Appel Maths For Physicists PDF

Mathematics for Physics
and Physicists
Mathematics for Physics

and Physicists
Walter Appel
Translated by Emmanuel Kowalski
Princeton University Press

Princeton and Oxford
Copyright
2007 by Princeton University Press
Published by Princeton University Press, 41 William Street,

Princeton, New Jersey 08540
In the United Kingdom: Princeton University Press, 3 Market Place,
Woodstock, Oxfordshire OX20 1SY
Originally published in French under the title Mathmatiques pour la physique... et les
physiciens!, copyright
c 2001 by H&K Editions, Paris.
Translated from the French by Emmanuel Kowalski

Ouvrage publi avec le concours du Ministre franais charg de la culture Centre
national du livre.
[This book was published with the support of the French Ministry of Culture
Centre national du livre.]
All Rights Reserved

Library of Congress Control Number: 2006935871
ISBN-13: 978-0-691-13102-3
ISBN-10: 0-691-13102-3
British Library Cataloging-in-Publication Data is available
Printed on acid-free paper.
pup.princeton.edu
Printed in the United States of America
10 9 8 7 6 5 4 3 2 1
Formatted with LATEX 2 under a Linux environment.
Contents
A books apology
xviii
Index of notation
xxii
Reminders: convergence of sequences and series

1.1 The problem of limits in physics . . . . . . . . . . . . . . . . . . . . . .
1.1.a
Two paradoxes involving kinetic energy . . . . . . . . . . . . .
1.1.b Romeo, Juliet, and viscous fluids . . . . . . . . . . . . . . . . .
1.1.c
Potential wall in quantum mechanics . . . . . . . . . . . . . .
1.1.d Semi-infinite filter behaving as waveguide . . . . . . . . . . . .
1.2 Sequences . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
1.2.a Sequences in a normed vector space . . . . . . . . . . . . . . .
1.2.b Cauchy sequences . . . . . . . . . . . . . . . . . . . . . . . . . .
1.2.c The fixed point theorem . . . . . . . . . . . . . . . . . . . . . .
1.2.d Double sequences . . . . . . . . . . . . . . . . . . . . . . . . . .
1.2.e Sequential definition of the limit of a function . . . . . . . .
1.2.f
Sequences of functions . . . . . . . . . . . . . . . . . . . . . . .
1.3 Series . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
1.3.a Series in a normed vector space . . . . . . . . . . . . . . . . .
1.3.b Doubly infinite series . . . . . . . . . . . . . . . . . . . . . . .
1.3.c Convergence of a double series . . . . . . . . . . . . . . . . . .
1.3.d Conditionally convergent series, absolutely convergent series .
1.3.e Series of functions . . . . . . . . . . . . . . . . . . . . . . . . .
1.4 Power series, analytic functions . . . . . . . . . . . . . . . . . . . . . . .
1.4.a Taylor formulas . . . . . . . . . . . . . . . . . . . . . . . . . . .
1.4.b Some numerical illustrations . . . . . . . . . . . . . . . . . . .
1.4.c Radius of convergence of a power series . . . . . . . . . . . . .
1.4.d Analytic functions . . . . . . . . . . . . . . . . . . . . . . . . .
1.5 A quick look at asymptotic and divergent series . . . . . . . . . . . . .
1.5.a Asymptotic series . . . . . . . . . . . . . . . . . . . . . . . . . .
1.5.b Divergent series and asymptotic expansions . . . . . . . . . .
Exercises . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
Problem . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
Solutions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
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47
Measure theory and the Lebesgue integral

2.1 The integral according to Mr. Riemann . . . .
2.1.a Riemann sums . . . . . . . . . . . . .
2.1.b Limitations of Riemanns definition .
2.2 The integral according to Mr. Lebesgue . . . .
2.2.a Principle of the method . . . . . . . .
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vi
Contents
2.2.b Borel subsets . . . . . . . . . . . . . . . . . .
2.2.c Lebesgue measure . . . . . . . . . . . . . . . .
2.2.d The Lebesgue -algebra . . . . . . . . . . . .
2.2.e Negligible sets . . . . . . . . . . . . . . . . . .
2.2.f
Lebesgue measure on Rn . . . . . . . . . . .
2.2.g Definition of the Lebesgue integral . . . . .
2.2.h Functions zero almost everywhere, space L1
2.2.i
And today? . . . . . . . . . . . . . . . . . . .
Exercises . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
Solutions . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
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Integral calculus
3.1 Integrability in practice . . . . . . . . . . . . . . . . . . .
3.1.a Standard functions . . . . . . . . . . . . . . . . .
3.1.b Comparison theorems . . . . . . . . . . . . . . .
3.2 Exchanging integrals and limits or series . . . . . . . . .
3.3 Integrals with parameters . . . . . . . . . . . . . . . . . .
3.3.a Continuity of functions defined by integrals . .
3.3.b Differentiating under the integral sign . . . . .
3.3.c Case of parameters appearing in the integration
3.4 Double and multiple integrals . . . . . . . . . . . . . . .
3.5 Change of variables . . . . . . . . . . . . . . . . . . . . .
Exercises . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
Solutions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
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116
Complex Analysis I
4.1 Holomorphic functions . . . . . . . . . . . . . . . . . . . . . .
4.1.a Definitions . . . . . . . . . . . . . . . . . . . . . . . .
4.1.b Examples . . . . . . . . . . . . . . . . . . . . . . . . .
4.1.c The operators / z and /z . . . . . . . . . . . . .
4.2 Cauchys theorem . . . . . . . . . . . . . . . . . . . . . . . . .
4.2.a Path integration . . . . . . . . . . . . . . . . . . . . . .
4.2.b Integrals along a circle . . . . . . . . . . . . . . . . . .
4.2.c Winding number . . . . . . . . . . . . . . . . . . . . .
4.2.d Various forms of Cauchys theorem . . . . . . . . . .
4.2.e Application . . . . . . . . . . . . . . . . . . . . . . . .
4.3 Properties of holomorphic functions . . . . . . . . . . . . . .
4.3.a The Cauchy formula and applications . . . . . . . . .
4.3.b Maximum modulus principle . . . . . . . . . . . . . .
4.3.c Other theorems . . . . . . . . . . . . . . . . . . . . . .
4.3.d Classification of zero sets of holomorphic functions
4.4 Singularities of a function . . . . . . . . . . . . . . . . . . . . .
4.4.a Classification of singularities . . . . . . . . . . . . . .
4.4.b Meromorphic functions . . . . . . . . . . . . . . . . .
4.5 Laurent series . . . . . . . . . . . . . . . . . . . . . . . . . . . .
4.5.a Introduction and definition . . . . . . . . . . . . . .
4.5.b Examples of Laurent series . . . . . . . . . . . . . . .
4.5.c The Residue theorem . . . . . . . . . . . . . . . . . . .
4.5.d Practical computations of residues . . . . . . . . . . .
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Contents
vii
4.6
Applications to the computation of horrifying integrals or ghastly sums 117

4.6.a Jordans lemmas . . . . . . . . . . . . . . . . . . . . . . . . . . 117
4.6.b Integrals on R of a rational function . . . . . . . . . . . . . . 118
4.6.c Fourier integrals . . . . . . . . . . . . . . . . . . . . . . . . . . 120
4.6.d Integral on the unit circle of a rational function . . . . . . . 121
4.6.e Computation of infinite sums . . . . . . . . . . . . . . . . . . 122
Exercises . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 125
Problem . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 128
Solutions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 129
Complex Analysis II
5.1 Complex logarithm; multivalued functions . . . .
5.1.a The complex logarithms . . . . . . . . . .
5.1.b The square root function . . . . . . . . .
5.1.c Multivalued functions, Riemann surfaces
5.2 Harmonic functions . . . . . . . . . . . . . . . . .
5.2.a Definitions . . . . . . . . . . . . . . . . .
5.2.b Properties . . . . . . . . . . . . . . . . . .
5.2.c A trick to find f knowing u . . . . . . .
5.3 Analytic continuation . . . . . . . . . . . . . . . .
5.4 Singularities at infinity . . . . . . . . . . . . . . .
5.5 The saddle point method . . . . . . . . . . . . . .
5.5.a The general saddle point method . . . .
5.5.b The real saddle point method . . . . . . .
Exercises . . . . . . . . . . . . . . . . . . . . . . . . . . .
Solutions . . . . . . . . . . . . . . . . . . . . . . . . . . .
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Conformal maps
6.1 Conformal maps . . . . . . . . . . . . . . . . . . . . . . . .
6.1.a Preliminaries . . . . . . . . . . . . . . . . . . . . .
6.1.b The Riemann mapping theorem . . . . . . . . . .
6.1.c Examples of conformal maps . . . . . . . . . . . .
6.1.d The Schwarz-Christoffel transformation . . . . . .
6.2 Applications to potential theory . . . . . . . . . . . . . . .
6.2.a Application to electrostatics . . . . . . . . . . . . .
6.2.b Application to hydrodynamics . . . . . . . . . . .
6.2.c Potential theory, lightning rods, and percolation
6.3 Dirichlet problem and Poisson kernel . . . . . . . . . . . .
Exercises . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
Solutions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
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155
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Distributions I
7.1 Physical approach . . . . . . . . . . . . . . . . . . . . . . . .
7.1.a
The problem of distribution of charge . . . . . .
7.1.b The problem of momentum and forces during an
7.2 Definitions and examples of distributions . . . . . . . . .
7.2.a Regular distributions . . . . . . . . . . . . . . . . .
7.2.b Singular distributions . . . . . . . . . . . . . . . .
7.2.c
Support of a distribution . . . . . . . . . . . . . .
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179
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185
187
viii
7.2.d Other examples . . . . . . . . . . . . . . . . . . . . . . . . . . .

Elementary properties. Operations . . . . . . . . . . . . . . . . . . . . .
7.3.a Operations on distributions . . . . . . . . . . . . . . . . . . .
7.3.b Derivative of a distribution . . . . . . . . . . . . . . . . . . . .
Dirac and its derivatives . . . . . . . . . . . . . . . . . . . . . . . . . . .
7.4.a The Heaviside distribution . . . . . . . . . . . . . . . . . . . .
7.4.b Multidimensional Dirac distributions . . . . . . . . . . . . . .
7.4.c
The distribution . . . . . . . . . . . . . . . . . . . . . . . . .
7.4.d Composition of with a function . . . . . . . . . . . . . . .
7.4.e
Charge and current densities . . . . . . . . . . . . . . . . . . .
Derivation of a discontinuous function . . . . . . . . . . . . . . . . . .
7.5.a Derivation of a function discontinuous at a point . . . . . .
7.5.b Derivative of a function with discontinuity along a surface S
7.5.c
Laplacian of a function discontinuous along a surface S . .
7.5.d Application: laplacian of 1/r in 3-space . . . . . . . . . . . . .
Convolution . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
7.6.a The tensor product of two functions . . . . . . . . . . . . . .
7.6.b The tensor product of distributions . . . . . . . . . . . . . . .
7.6.c
Convolution of two functions . . . . . . . . . . . . . . . . . .
7.6.d Fuzzy measurement . . . . . . . . . . . . . . . . . . . . . . .
7.6.e
Convolution of distributions . . . . . . . . . . . . . . . . . . .
7.6.f
Applications . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
7.6.g The Poisson equation . . . . . . . . . . . . . . . . . . . . . . .
Physical interpretation of convolution operators . . . . . . . . . . . . .
Discrete convolution . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
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216
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220
Distributions II
8.1 Cauchy principal value . . . . . . . . . . . . . . . . . . . . . . . . . . .
8.1.a Definition . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
8.1.b Application to the computation of certain integrals . . . . . .
8.1.c Feynmans notation . . . . . . . . . . . . . . . . . . . . . . . .
8.1.d Kramers-Kronig relations . . . . . . . . . . . . . . . . . . . . .
8.1.e
A few equations in the sense of distributions . . . . . . . . . .
8.2 Topology in D . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
8.2.a Weak convergence in D . . . . . . . . . . . . . . . . . . . . . .
8.2.b Sequences of functions converging to . . . . . . . . . . . . .
8.2.c Convergence in D and convergence in the sense of functions
8.2.d Regularization of a distribution . . . . . . . . . . . . . . . . .
8.2.e Continuity of convolution . . . . . . . . . . . . . . . . . . . .
8.3 Convolution algebras . . . . . . . . . . . . . . . . . . . . . . . . . . . .
8.4 Solving a differential equation with initial conditions . . . . . . . . .
8.4.a First order equations . . . . . . . . . . . . . . . . . . . . . . . .
8.4.b The case of the harmonic oscillator . . . . . . . . . . . . . . .
8.4.c Other equations of physical origin . . . . . . . . . . . . . . . .
Exercises . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
Problem . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
Solutions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
223
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230
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236
238
238
239
240
241
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7.3
7.4
7.5
7.6
7.7
7.8
8
Contents
Contents
9
ix
Hilbert spaces; Fourier series

9.1 Insufficiency of vector spaces . . . . . . . . . . . . . .
9.2 Pre-Hilbert spaces . . . . . . . . . . . . . . . . . . . . .
9.2.a The finite-dimensional case . . . . . . . . . .
9.2.b Projection on a finite-dimensional subspace
9.2.c Bessel inequality . . . . . . . . . . . . . . . .
9.3 Hilbert spaces . . . . . . . . . . . . . . . . . . . . . . .
9.3.a Hilbert basis . . . . . . . . . . . . . . . . . . .
9.3.b The 2 space . . . . . . . . . . . . . . . . . . .
9.3.c The space L2 [0, a] . . . . . . . . . . . . . . .
9.3.d The L2 (R) space . . . . . . . . . . . . . . . .
9.4 Fourier series expansion . . . . . . . . . . . . . . . . .
9.4.a Fourier coefficients of a function . . . . . .
9.4.b Mean-square convergence . . . . . . . . . . .
9.4.c Fourier series of a function f L1 [0, a] . .
9.4.d Pointwise convergence of the Fourier series .
9.4.e Uniform convergence of the Fourier series .
9.4.f
The Gibbs phenomenon . . . . . . . . . . . .
Exercises . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
Problem . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
Solutions . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
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277
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296
11 Fourier transform of distributions

11.1 Definition and properties . . . . . . . . . . . . . . .
11.1.a Tempered distributions . . . . . . . . . . . .
11.1.b Fourier transform of tempered distributions
11.1.c Examples . . . . . . . . . . . . . . . . . . . .
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299
299
300
301
303
10 Fourier transform of functions

10.1 Fourier transform of a function in L1 . . .
10.1.a Definition . . . . . . . . . . . . . . .
10.1.b Examples . . . . . . . . . . . . . . .
10.1.c The L1 space . . . . . . . . . . . . .
10.1.d Elementary properties . . . . . . . .
10.1.e Inversion . . . . . . . . . . . . . . . .
10.1.f Extension of the inversion formula
10.2 Properties of the Fourier transform . . . . .
10.2.a Transpose and translates . . . . . . .
10.2.b Dilation . . . . . . . . . . . . . . . .
10.2.c Derivation . . . . . . . . . . . . . . .
10.2.d Rapidly decaying functions . . . . .
10.3 Fourier transform of a function in L2 . . .
10.3.a The space S . . . . . . . . . . . . .
10.3.b The Fourier transform in L2 . . . .
10.4 Fourier transform and convolution . . . . .
10.4.a Convolution formula . . . . . . . .
10.4.b Cases of the convolution formula .
Exercises . . . . . . . . . . . . . . . . . . . . . . . .
Solutions . . . . . . . . . . . . . . . . . . . . . . . .
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Contents
11.1.d Higher-dimensional Fourier transforms . . . . . . . . . . . .
11.1.e Inversion formula . . . . . . . . . . . . . . . . . . . . . . . .
11.2 The Dirac comb . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
11.2.a Definition and properties . . . . . . . . . . . . . . . . . . . .
11.2.b Fourier transform of a periodic function . . . . . . . . . . .
11.2.c Poisson summation formula . . . . . . . . . . . . . . . . . .
11.2.d Application to the computation of series . . . . . . . . . . .
11.3 The Gibbs phenomenon . . . . . . . . . . . . . . . . . . . . . . . . .
11.4 Application to physical optics . . . . . . . . . . . . . . . . . . . . . .
11.4.a Link between diaphragm and diffraction figure . . . . . . .
11.4.b Diaphragm made of infinitely many infinitely narrow slits
11.4.c Finite number of infinitely narrow slits . . . . . . . . . . . .
11.4.d Finitely many slits with finite width . . . . . . . . . . . . . .
11.4.e Circular lens . . . . . . . . . . . . . . . . . . . . . . . . . . .
11.5 Limitations of Fourier analysis and wavelets . . . . . . . . . . . . . .
Exercises . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
Problem . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
Solutions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
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305
306
307
307
308
309
310
311
314
314
315
316
318
320
321
324
325
326
12 The Laplace transform

12.1 Definition and integrability . . . . . . . . . . . . . . . . . . . .
12.1.a Definition . . . . . . . . . . . . . . . . . . . . . . . . . .
12.1.b Integrability . . . . . . . . . . . . . . . . . . . . . . . . .
12.1.c Properties of the Laplace transform . . . . . . . . . . .
12.2 Inversion . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
12.3 Elementary properties and examples of Laplace transforms . .
12.3.a Translation . . . . . . . . . . . . . . . . . . . . . . . . .
12.3.b Convolution . . . . . . . . . . . . . . . . . . . . . . . .
12.3.c Differentiation and integration . . . . . . . . . . . . . .
12.3.d Examples . . . . . . . . . . . . . . . . . . . . . . . . . .
12.4 Laplace transform of distributions . . . . . . . . . . . . . . . .
12.4.a Definition . . . . . . . . . . . . . . . . . . . . . . . . . .
12.4.b Properties . . . . . . . . . . . . . . . . . . . . . . . . . .
12.4.c Examples . . . . . . . . . . . . . . . . . . . . . . . . . .
12.4.d The z-transform . . . . . . . . . . . . . . . . . . . . . . .
12.4.e Relation between Laplace and Fourier transforms . . .
12.5 Physical applications, the Cauchy problem . . . . . . . . . . .
12.5.a Importance of the Cauchy problem . . . . . . . . . . .
12.5.b A simple example . . . . . . . . . . . . . . . . . . . . . .
12.5.c Dynamics of the electromagnetic field without sources
Exercises . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
Solutions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
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336
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339
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341
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342
342
344
344
345
346
346
347
348
351
352
13 Physical applications of the Fourier transform

13.1 Justification of sinusoidal regime analysis . . . . . .
13.2 Fourier transform of vector fields: longitudinal and
13.3 Heisenberg uncertainty relations . . . . . . . . . . .
13.4 Analytic signals . . . . . . . . . . . . . . . . . . . . .
13.5 Autocorrelation of a finite energy function . . . . .
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355
355
358
359
365
368
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transverse
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Contents
13.5.a Definition . . . . . . . . . . . . . . . . . . . . . .
13.5.b Properties . . . . . . . . . . . . . . . . . . . . . .
13.5.c Intercorrelation . . . . . . . . . . . . . . . . . . .
13.6 Finite power functions . . . . . . . . . . . . . . . . . . .
13.6.a Definitions . . . . . . . . . . . . . . . . . . . . .
13.6.b Autocorrelation . . . . . . . . . . . . . . . . . . .
13.7 Application to optics: the Wiener-Khintchine theorem
Exercises . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
Solutions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
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368
368
369
370
370
370
371
375
376
14 Bras, kets, and all that sort of thing

14.1 Reminders about finite dimension . . . . . . . . . . . . . . . . .
14.1.a Scalar product and representation theorem . . . . . . . .
14.1.b Adjoint . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
14.1.c Symmetric and hermitian endomorphisms . . . . . . . .
14.2 Kets and bras . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
14.2.a Kets | H . . . . . . . . . . . . . . . . . . . . . . . . .
14.2.b Bras | H . . . . . . . . . . . . . . . . . . . . . . . . .
14.2.c Generalized bras . . . . . . . . . . . . . . . . . . . . . . .
14.2.d Generalized
kets . . . . . . . . . . . . . . . . . . . . . . .
P
14.2.e Id = n |n n | . . . . . . . . . . . . . . . . . . . . . .
14.2.f Generalized basis . . . . . . . . . . . . . . . . . . . . . . .
14.3 Linear operators . . . . . . . . . . . . . . . . . . . . . . . . . . . .
14.3.a Operators . . . . . . . . . . . . . . . . . . . . . . . . . . .
14.3.b Adjoint . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
14.3.c Bounded operators, closed operators, closable operators
14.3.d Discrete and continuous spectra . . . . . . . . . . . . . .
14.4 Hermitian operators; self-adjoint operators . . . . . . . . . . . .
14.4.a Definitions . . . . . . . . . . . . . . . . . . . . . . . . . .
14.4.b Eigenvectors . . . . . . . . . . . . . . . . . . . . . . . . . .
14.4.c Generalized eigenvectors . . . . . . . . . . . . . . . . . . .
14.4.d Matrix representation . . . . . . . . . . . . . . . . . . .
14.4.e Summary of properties of the operators P and X . . . .
Exercises . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
Solutions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
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382
383
384
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387
387
389
390
391
393
394
396
397
398
401
403
404
15 Green functions
407
15.1 Generalities about Green functions . . . . . . . . . . . . . . . . . . . . 407
15.2 A pedagogical example: the harmonic oscillator . . . . . . . . . . . . 409
15.2.a Using the Laplace transform . . . . . . . . . . . . . . . . . . . 410
15.2.b Using the Fourier transform . . . . . . . . . . . . . . . . . . . 410
15.3 Electromagnetism and the dAlembertian operator . . . . . . . . . . . 414
15.3.a Computation of the advanced and retarded Green functions 414
15.3.b Retarded potentials . . . . . . . . . . . . . . . . . . . . . . . . . 418
15.3.c Covariant expression of advanced and retarded Green functions 421
15.3.d Radiation . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 421
15.4 The heat equation . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 422
15.4.a One-dimensional case . . . . . . . . . . . . . . . . . . . . . . . 423
15.4.b Three-dimensional case . . . . . . . . . . . . . . . . . . . . . . 426
xii
Contents
15.5 Quantum mechanics . . . . . . . . . . . . . . . . . . . . . . . . . . . . 427

15.6 Klein-Gordon equation . . . . . . . . . . . . . . . . . . . . . . . . . . . 429
Exercises . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 432
16 Tensors
16.1 Tensors in affine space . . . . . . . . . . . . . . . . . . . . . . .
16.1.a Vectors . . . . . . . . . . . . . . . . . . . . . . . . . . . .
16.1.b Einstein convention . . . . . . . . . . . . . . . . . . . .
16.1.c Linear forms . . . . . . . . . . . . . . . . . . . . . . . .
16.1.d Linear maps . . . . . . . . . . . . . . . . . . . . . . . . .
16.1.e Lorentz transformations . . . . . . . . . . . . . . . . . .
16.2 Tensor product of vector spaces: tensors . . . . . . . . . . . . .
16.2.a Existence of the tensor product of two vector spaces .
0
16.2.b Tensor product of linear forms: tensors of type 2 . .
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16.2.c Tensor product of vectors: tensors of type 0 . . . . .
16.2.d Tensor product of a vector and a linear form: linear
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16.2.e Tensors of type q . . . . . . . . . . . . . . . . . . . . .
16.3 The metric, or, how to raise and lower indices . . . . . . . . .
16.3.a Metric and pseudo-metric . . . . . . . . . . . . . . . . .
16.3.b Natural duality by means of the metric . . . . . . . . .
16.3.c Gymnastics: raising and lowering indices . . . . . . . .
16.4 Operations on tensors . . . . . . . . . . . . . . . . . . . . . . .
16.5 Change of coordinates . . . . . . . . . . . . . . . . . . . . . . .
16.5.a Curvilinear coordinates . . . . . . . . . . . . . . . . . .
16.5.b Basis vectors . . . . . . . . . . . . . . . . . . . . . . . . .
16.5.c Transformation of physical quantities . . . . . . . . . .
16.5.d Transformation of linear forms . . . . . . . . . . . . . .
16.5.e Transformation of an arbitrary tensor field . . . . . . .
16.5.f Conclusion . . . . . . . . . . . . . . . . . . . . . . . . .
Solutions . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
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458
459
460
461
462
17 Differential forms
17.1 Exterior algebra . . . . . . . . . . . . . . . . . . . . . . . .
17.1.a 1-forms . . . . . . . . . . . . . . . . . . . . . . . . .
17.1.b Exterior 2-forms . . . . . . . . . . . . . . . . . . .
17.1.c Exterior k-forms . . . . . . . . . . . . . . . . . . .
17.1.d Exterior product . . . . . . . . . . . . . . . . . . .
17.2 Differential forms on a vector space . . . . . . . . . . . .
17.2.a Definition . . . . . . . . . . . . . . . . . . . . . . .
17.2.b Exterior derivative . . . . . . . . . . . . . . . . . .
17.3 Integration of differential forms . . . . . . . . . . . . . . .
17.4 Poincars theorem . . . . . . . . . . . . . . . . . . . . . .
17.5 Relations with vector calculus: gradient, divergence, curl
17.5.a Differential forms in dimension 3 . . . . . . . . .
17.5.b Existence of the scalar electrostatic potential . . .
17.5.c Existence of the vector potential . . . . . . . . . .
17.5.d Magnetic monopoles . . . . . . . . . . . . . . . . .
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433
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441
443
Contents
xiii
17.6 Electromagnetism in the language of differential forms . . . . . . . . 480

Problem . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 484
Solution . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 485
18 Groups and group representations
18.1 Groups . . . . . . . . . . . . . . .
18.2 Linear representations of groups
18.3 Vectors and the group SO(3) . .
18.4 The group SU(2) and spinors . .
18.5 Spin and Riemann sphere . . . .
Exercises . . . . . . . . . . . . . . . . . .
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489
489
491
492
497
503
505
19 Introduction to probability theory

19.1 Introduction . . . . . . . . . .
19.2 Basic definitions . . . . . . . .
19.3 Poincar formula . . . . . . .
19.4 Conditional probability . . .
19.5 Independent events . . . . . .
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509
510
512
516
517
519
20 Random variables
20.1 Random variables and probability distributions . . . .
20.2 Distribution function and probability density . . . . .
20.2.a Discrete random variables . . . . . . . . . . . . .
20.2.b (Absolutely) continuous random variables . . .
20.3 Expectation and variance . . . . . . . . . . . . . . . . .
20.3.a Case of a discrete r.v. . . . . . . . . . . . . . . . .
20.3.b Case of a continuous r.v. . . . . . . . . . . . . .
20.4 An example: the Poisson distribution . . . . . . . . . .
20.4.a Particles in a confined gas . . . . . . . . . . . . .
20.4.b Radioactive decay . . . . . . . . . . . . . . . . . .
20.5 Moments of a random variable . . . . . . . . . . . . . .
20.6 Random vectors . . . . . . . . . . . . . . . . . . . . . . .
20.6.a Pair of random variables . . . . . . . . . . . . . .
20.6.b Independent random variables . . . . . . . . . .
20.6.c Random vectors . . . . . . . . . . . . . . . . . .
20.7 Image measures . . . . . . . . . . . . . . . . . . . . . . .
20.7.a Case of a single random variable . . . . . . . . .
20.7.b Case of a random vector . . . . . . . . . . . . . .
20.8 Expectation and characteristic function . . . . . . . . .
20.8.a Expectation of a function of random variables
20.8.b Moments, variance . . . . . . . . . . . . . . . . .
20.8.c Characteristic function . . . . . . . . . . . . . .
20.8.d Generating function . . . . . . . . . . . . . . . .
20.9 Sum and product of random variables . . . . . . . . . .
20.9.a Sum of random variables . . . . . . . . . . . . .
20.9.b Product of random variables . . . . . . . . . . .
20.9.c Example: Poisson distribution . . . . . . . . . .
20.10 Bienaym-Tchebychev inequality . . . . . . . . . . . . .
20.10.a Statement . . . . . . . . . . . . . . . . . . . . . .
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521
521
524
526
526
527
527
528
530
530
531
532
534
534
537
538
539
539
540
540
540
541
541
543
543
543
546
547
547
547
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xiv
Contents
20.10.b Application: Buffons needle . . . . . . . . . . . . . . . . . . . 549

20.11 Independance, correlation, causality . . . . . . . . . . . . . . . . . . . 550
21 Convergence of random variables: central limit
21.1 Various types of convergence . . . . . . . . .
21.2 The law of large numbers . . . . . . . . . .
21.3 Central limit theorem . . . . . . . . . . . .
Exercises . . . . . . . . . . . . . . . . . . . . . . . .
Problems . . . . . . . . . . . . . . . . . . . . . . . .
Solutions . . . . . . . . . . . . . . . . . . . . . . . .
theorem
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553
553
555
556
560
563
564
A Reminders concerning topology and normed vector spaces

A.1 Topology, topological spaces . . . . . . . . . . . . . . . .
A.2 Normed vector spaces . . . . . . . . . . . . . . . . . . . .
A.2.a Norms, seminorms . . . . . . . . . . . . . . . . .
A.2.b Balls and topology associated to the distance . .
A.2.c Comparison of sequences . . . . . . . . . . . . .
A.2.d Bolzano-Weierstrass theorems . . . . . . . . . . .
A.2.e Comparison of norms . . . . . . . . . . . . . . .
A.2.f Norm of a linear map . . . . . . . . . . . . . . .
Exercise . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
Solution . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . .
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573
573
577
577
578
580
581
581
583
583
584
B Elementary reminders of differential calculus

B.1 Differential of a real-valued function . . . . . .
B.1.a Functions of one real variable . . . . .
B.1.b Differential of a function f : Rn R
B.1.c Tensor notation . . . . . . . . . . . . . .
B.2 Differential of map with values in R p . . . . .
B.3 Lagrange multipliers . . . . . . . . . . . . . . . .
Solution . . . . . . . . . . . . . . . . . . . . . . . . . .
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585
585
585
586
587
587
588
591
C Matrices
C.1 Duality . . . . . . . . . . . . . . . . . . . . . . .
C.2 Application to matrix representation . . . . . .
C.2.a Matrix representing a family of vectors
C.2.b Matrix of a linear map . . . . . . . . .
C.2.c Change of basis . . . . . . . . . . . . . .
C.2.d Change of basis formula . . . . . . . .
C.2.e Case of an orthonormal basis . . . . .
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593
593
594
594
594
595
595
596
Appendices
D A few proofs
597
Contents
xv
Tables
Fourier transforms
609
Laplace transforms
613
Probability laws
616
Further reading
617
References
621
Portraits
627
Sidebars
629
Index
631
To Kirone Mallick, Angel Alastuey and Michael Kiessling,

for their friendship and our discussions scientific or otherwise.
Thanks
Many are those who I wish to thank for this book: Jean-Franois Colombeau
who was in charge of the course of Mathematics for the Magistre des sciences de la
matire of the ENS Lyon before me; Michel Peyrard who then asked me to replace
him; Vronique Terras and Jean Farago who were in charge of the exercise sessions
during the three years when I taught this course.
Many friends and colleagues were kind enough to read and re-read one or more
chapters: Julien Barr, Maxime Clusel, Thierry Dauxois, Kirone Mallick, Julien
Michel, Antoine Naert, Catherine Ppin, Magali Ribaut, Erwan Saint-Loubert Bi;
and especially Paul Pichaureau, TEX-guru, who gave many (very) sharp and (always)
pertinent comments concerning typography and presentation as well as contents. I
also wish to thank the ditions H&K and in particular Sbastien Desreux, for his
hard work and stimulating demands; Brangre Condomines for her attentive reading.
I am also indebted to Jean-Franois Quint for many long and fascinating mathematical discussions.
This edition owes a lot to Emmanuel Kowalski, friend of many years, Professor at
the University of Bordeaux, who led me to refine certain statements, correct others,
and more generally helped to clarify some delicate points.
For sundry diverse reasons, no less important, I want to thank Craig Thompson,
author of Good-bye Chunky Rice [89]and Blankets [90], who created the illustrations for
Chapter 1; Jean-Paul Marchal, master typographer and maker of popular images in
pinal; Angel Alastuey who taught me so much in physics and mathematics; Frdric,
Latitia, Samuel, Koupaa and Alanis Vivien for their support and friendship; Claude
Garcia, for his joie de vivre, his fideos, his fine bottles, his knowledge of topology, his
advice on Life, the Universe... and the Rest! without forgetting Anita, Alice and Hugo
for their affection.
I must not forget the many readers who, since the first French edition, have
communicated their remarks and corrections; in particular Jean-Julien Fleck, Marc
Rezzouk, Franoise Cornu, Cline Chevalier and professors Jean Cousteix, Andreas
de Vries, and Franois Thirioux; without forgetting all the students of the Magistre
des sciences de la matire of ENS Lyon, promotions 1993 to 1999, who followed
my classes and exercise sessions and who contributed, through their remarks and
questions, to the elaboration and maturation of this book.
Et, bien sr, les plus importants, Anne-Julia, Solveig et Anton avec tout mon
amour.
Although the text has been very consciously written, read, and proofread, errors, omissions or imprecisions may still remain. The author welcomes any remark, criticism, or correction that a reader may wish to
communicate, care of the publisher, for instance, by means of the email address Errare.humanum.est@h-k.fr.
A list of errata will be kept updated on the web site
http://press.princeton.edu
A books apology
hy should a physicist study mathematics? There is a fairly fashionable current of thought that holds that the use of advanced mathematics is of
little real use in physics, and goes sometimes as far as to say that knowing
mathematics is already by itself harmful to true physics. However, I am and remain
convinced that mathematics is still a precious source of insight, not only for students of
physics, but also for researchers.
Many only see mathematics as a tool and of course, it is in part a tool, but
they should be reminded that, as Galileo said, the book of Nature is written in the
mathematicians language.1 Since Galileo and Newton, the greatest physicists give
examples that knowing mathematics provides the means to understand precise physical
notions, to use them more easily, to establish them on a sure foundation, and even
more importantly, to discover new ones.2 In addition to ensuring a certain rigor in
reasoning which is indispensable in any scientific study, mathematics belongs to the
natural lexicon of physicists. Even if the rules of proportion and the fundamentals of
calculus are sufficient for most purposes, it is clear that a richer vocabulary can lead
to much deeper insights. Imagine if Shakespeare or Joyce had only had a few hundred
words to chose from!
Is is therefore discouraging to sometimes hear physicists dismiss certain theories because this is only mathematics. In fact, the two disciplines are so closely related that
the most prestigious mathematical award, the Fields medal, was given to the physicist
Edward Witten in 1990, rewarding him for the remarkable mathematical discoveries
that his ideas led to.
How should you read this book? Or rather how should you not read some parts?
1
Philosophy is written is this vast book which is forever open before our eyes I mean, the Universe
but it can not be read until you have learned the tongue and understood the character in which it is
written. It is written in the language of mathematics, and its characters are triangles, circles, and other
geometric pictures, without which it is not humanly possible to understand a single word [...].
2
I will only mention Newton (gravitation, differential and integral calculus), Gauss (optics,
magnetism, all the mathematics of his time, and quite a bit that was only understood much
later), Hamilton (mechanics, differential equations, algebra), Heaviside (symbolic calculus,
signal theory), Gibbs (thermodynamics, vector analysis) and of course Einstein. One could
write a much longer list. If Richard Feynman presents a very physical description of his art
in his excellent physics course [35], appealing to remarkably little formalism, it is nevertheless
the fact that he was the master of the elaborate mathematics involved, as his research works
show.
A books apology
xix
Since the reader may want to learn first certain specific topics among those present, here
is a short description of the contents of this text:
The first chapter serves as a reminder of some simple facts concerning the fundamental notion of convergence. There shouldnt be much in the way of new
mathematics there for anyone who has followed a rigorous calculus course,
but there are many counter-examples to be aware of. Most importantly, a long
section describes the traps and difficulties inherent in the process of exchanging
two limits in the setting of physical models. It is not always obvious where,
in physical reasoning, one has to exchange two mathematical limits, and many
apparent paradoxes follow from this fact.
The real beginning concerns the theory of integration, which is briefly presented
in the form of the Lebesgue integral based on measure theory (Chapter 2). For
many readers, this may be omitted in a first reading. Chapter 3 discusses the
basic results and techniques of integral calculus.
Chapters 4 to 6 present the theory of functions of a complex variable, with a
number of applications:
the residue method, which is an amazing tool for integral calculus;
some physical notions, such as causality, are very closely related to analyticity of functions on the complex plane (Section 13.4);
harmonic functions (such that f = 0) in two dimensions are linked to
the real parts of holomorphic (analytic) functions (Chapter 5);
conformal maps (those which preserve angles) can be used to simplify
boundary conditions in problems of hydrodynamics or electromagnetism
(Chapter 6);
Chapters 7 and 8 concern the theory of distributions (generalized functions)
and their applications to physics. These form a relatively self-contained subset of
the book.
Chapters 9 to 12 deal with Hilbert spaces, Fourier series, Fourier and Laplace
transforms, which have too many physical applications to attempt a list. Chapter 13 develops some of those applications, and this chapter also requires complex
analysis.
Chapter 14 is a short (probably too short) introduction to the Dirac notation
used in quantum mechanics: kets | and bras |. The notions of generalized
eigenbasis and self-adjoint operators on Hilbert space are also discussed.
Several precise physical problems are considered and solved in Chapter 15 by the
method of Green functions. This method is usually omitted from textbooks on
electromagnetism (where a solution is taken straight out of a magicians hat) or
of field theory (where it is assumed that the method is known). I hope to fill
a gap for students by presenting the necessary (and fairly simple) computations
from beginning to end, using physicists notation.
xx
A books apology
Chapters 16 and 17 about tensor calculus and differential forms are also somewhat independent from the rest of the text. Those two chapters are only brief
introductions to their respective topics.
Chapter 18 has the modest goal of relating some notions of topology and group
theory to the idea of spin in quantum mechanics.
Probability theory, discussed in Chapters 19 to 21, is almost entirely absent
from the standard physics curriculum, although the basic vocabulary and results of probability seem necessary to any physicist interested in theory (stochastic
equations, Brownian motion, quantum mechanics and statistical mechanics all
require probability theory) or experiments (Gaussian white noise, measurement
errors, standard deviation of a data set...)
Finally, a few appendices contain further reminders of elementary mathematical
notions and the proofs of some interesting results, the length of which made their
insertion in the main text problematic.
Many physical applications, using mathematical tools with the usual notation of physics, are included in the text. They can be found by looking in the
index at the items under Physical applications.
A books apology
xxi
Remarks on the biographical summaries

The short biographies of mathematicians which are interspersed in the text are taken from
many sources:
Bertrand Hauchecorne and Daniel Suratteau, Des mathmaticiens de A Z, Ellipses,
1996 (French).
The web site of Saint-Andrews University (Scotland)
www-history.mcs.st-andrews.ac.uk/history/Mathematicians
and the web site of the University of Colorado at Boulder
www.colorado.edu/education/DMP
from which certain pictures are also taken.
The literary structure of scientific argument, edited by Peter Dear, University of Pennsylvania
Press, 1991.
Simon Gindikin, Histoires de mathmaticiens et de physiciens, Cassini, 2000.
Encyclopdia Universalis, Paris, 1990.
Translators foreword
I am a mathematician and have now forgotten most of the little physics I learned in school
(although Ive probably picked up a little bit again by translating this book). I would like to
mention here two more reasons to learn mathematics, and why this type of book is therefore
very important.
First, physicists benefit from knowing mathematics (in addition to the reasons Walter
mentioned) because, provided they immerse themselves in mathematics sufficiently to become
fluent in its language, they will gain access to new intuitions. Intuitions are very different from
any set of techniques, or tools, or methods, but they are just as indispensable for a researcher,
and they are the hardest to come by.3 A mathematicians intuitions are very different from
those of a physicist, and to have both available is an enormous advantage.
The second argument is different, and may be subjective: physics is hard, much harder
in some sense than mathematics. A very simple and fundamental physical problem may
be all but impossible to solve because of the complexity (apparent or real) of Nature. But
mathematicians know that a simple, well-formulated, natural mathematical problem (in some
sense that is impossible to quantify!) has a simple solution. This solution may require inventing
entirely new concepts, and may have to wait for a few hundred years before the idea comes to
a brilliant mind, but it is there. What this means is that if you manage to put the physical
problem in a very natural mathematical form, the guiding principles of mathematics may lead
you to the solution. Dirac was certainly a physicist with a strong sense of such possibilities;
this led him to discover antiparticles, for instance.
3
Often, nothing will let you understand the intuition behind some important idea except,
essentially, rediscovering by yourself the most crucial part of it.
Notation
Symbol
Meaning
,
(|)
|, |
X
1
def
=
[a, b]
]a, b[
[[1, n]]
without (example: A \ B)
has for Laplace transform
duality bracket: a, x = a(x)
scalar product
ket, bra
integral part
Dirac comb
constant function x 7 1
convolution product
tensor product
exterior product
boundary (example: , boundary of )
laplacian
dAlembertian
equivalence of sequences or functions
isomorphic to (example: E F )
approximate equality (in physics)
equality defining the left-hand side
closed interval : {x R ; a x b}
open interval : {x R ; a < x < b}
{1, . . . , n}
Page
333
251
379,380
186
211
210, 439
464
414
580
Index of notation
Latin letter
B
B(a ; r ), B (a ; r )
Bil(E F , G )
A
C
C (I , R)
C (a ; r )
Cnt
D
D
D+ , D
diag(1 , . . . , n )
e
( e )
E(X )
F [ f ], fe
fb
f
f (n)
GLn (K)
G(E )
H
K, K
2
L1
L2
L (E , F )
L
Mn (K)
P(A)
P(A)
R
Re(z), Im(z)
S
S
S, Sn
t
u, v, w, x
pv
X,Y,Z
Meaning
Borel -algebra
open (closed) ball centered at a with radius r
space of bilinear maps from E F to G
complement of A in
C {}
vector space of continuous real-valued functions on I
circle centered at a with radius r
any constant
vector space of test functions
vector space of distributions
vector space of distributions with bounded support
on the left (right)
square matrix with coefficients mi j = i i j
Nepers constant (e: electric charge)
basis of a vector space
expectation of the r.v. X
Fourier transform of f
Laplace transform of f
transpose of f ; f(x) = f (x)
n-th derivative of f
group of invertible matrices of order n over K
group of automorphisms of the vector space E
Heaviside distribution or function
R or C
P
space of sequences (un )nN such that |un |2 < +
space of integrable functions
space of square-integrable functions
space of linear maps from E to F
Lebesgue -algebra
algebra of square matrices of order n over K
set of subsets of A
probability of the event A
R {, +}
real or imaginary part of z
space of Schwartz functions
space of tempered distributions
group of permutations (of n elements)
transpose of a matrix : ( tA)i j = A j i
vectors in E
principal value
random variables
xxiii
Page
58
512
146
182
183
236
433
527
278
333
189
193
261
280
262
60
57
514
289
300
26
434
188
521
Index of notation
xxiv
Greek letter
(
)
,
, i j
()
k
(X )
A
,
dual basis of the basis (e )

Euler function
Dirac distribution, its derivative
Kronecker symbol
signature of the permutation
space of exterior k-forms
rectangle function
standard deviation of the r.v. X
characteristic function of A
differential forms
436
154
185, 196
445
465
464
213, 279
528
63
436
Abbreviations
(n).v.s.
r.v.
i.r.v.
(normed) vector space

random variable
independent random variables
521
537
Chapter
Reminders: convergence of
sequences and series
This first chapter, which is quite elementary, is essentially a survey of the notion
of convergence of sequences and series. Readers who are very confortable with this
concept may start reading the next chapter.
However, although the mathematical objects we discuss are well known in principle, they have some unexpected properties. We will see in particular that the order
of summation may be crucial to the evaluation of the series, so that changing the
order of summation may well change its sum.
We start this chapter by discussing two physical problems in which a limit process
is hidden. Each leads to an apparent paradox, which can only be resolved when the
underlying limit is explicitly brought to light.
1.1
The problem of limits in physics
1.1.a
Two paradoxes involving kinetic energy

First paradox
Consider a truck with mass m driving at constant speed v = 60 mph on a

perfectly straight stretch of highway (think of Montana). We assume that,
friction being negligible, the truck uses no gas to remain at constant speed.
Reminders concerning convergence
On the other hand, to accelerate, it must use (or guzzle) gallons of gas in
order to increase its kinetic energy by an amount of 1 joule. This assumption,
although it is imperfect, is physically acceptable because each gallon of gas
yields the same amount of energy.
So, when the driver decides to increase its speed to reach v = 80 mph, the
quantity of gas required to do so is equal to the difference of kinetic energy,
namely, it is
1
1
(E c E c ) = m(v 2 v 2 ) = m(6 400 3 600) = 1 400 m.
2
2
With m = 10 1000 J mile2 hour2 , say, this amounts to 0.14 gallon. Jolly
good.
Now, let us watch the same scene of the truck accelerating, as observed by
a highway patrolman, initially driving as fast as the truck w = v = 60 mph,
but with a motorcycle which is unable to go faster.
The patrolman, having his college physics classes at his fingertips, argues as
follows: in my own Galilean reference frame, the relative speed of the truck
was previously v = 0 and is now v = 20 mph. To do this, the amount of
gas it has guzzled is equal to the difference in kinetic energies:
1
1
(E c E c ) = m (v )2 (v )2 = m(400 0) = 200 m,
2
2
or around 0.02 gallons.
There is here a clear problem, and one of the two observers must be wrong.
Indeed, the Galilean relativity principle states that all Galilean reference frames
are equivalent, and computing kinetic energy in the patrolmans reference
frame is perfectly legitimate.
How is this paradox resolved?
We will come to the solution, but first here is another problem. The reader,
before going on to read the solutions, is earnestly invited to think and try to
solve the problem by herself.
Second paradox
Consider a highly elastic rubber ball in free fall as we first see it. At some
point, it hits the ground, and we assume that this is an elastic shock.
Most high-school level books will describe the following argument: assume that, at the instant t = 0 when the ball hits the ground, the speed of
the ball is v1 = 10 ms1 . Since the shock is elastic, there is conservation of
total energy before and after. Hence the speed of the ball after the rebound is
v2 = v1 , or simply +10 ms1 going up.
This looks convincing enough. But it is not so impressive if seen from
the point of view of an observer who is also moving down at constant speed
vobs = v1 = 10 ms1 . For this observer, the speed of the ball before the
shock is v1 = v1 vobs = 0 ms1 , so it has zero kinetic energy. However,
after rebounding, the speed of the ball is v2 = v2 vobs = 20 ms1 , and

therefore it has nonzero kinetic energy! With the analogue of the reasoning
above, one should still have found v2 = v1 = 0 (should the ball go through
the ground?)
So there is something fishy in this argument also. It is important to
remember that the fact that the right answer is found in the first case does not imply
that the argument that leads to the answer is itself correct.
Readers who have solved the first paradox will find no difficulty in this
second one.
Paradoxes resolved
Kinetic energy is of course not the same in every reference frame. But this
is not so much the kinetic energy we are interested in; rather, we want the
difference before and after the event described.
Lets go back to elementary mechanics. What happens, in two distinct
reference frames, to a system of N solid bodies with initial speed v i (i =
1, . . . , N ) and final speed vi after some shock?
In the first reference frame, the difference of kinetic energy is given by
E c =
N
X
i=1
m i ( v i 2 v2i ).
In a second reference frame, with relative speed w with respect to the first,
the difference is equal to
E c =
N
X
i=1
N
X
i=1

m i ( vi w)2 ( v i w)2
m i ( v i 2
vi 2) 2 w
N
X
i=1
m i ( v i v i )
= E c 2 w P ,
(we use as labels for any physical quantity expressed in the new reference
frame), so that E c = E c as long as the total momentum is preserved during
the shock, in other words if P = 0.
In the case of the truck and the patrolman, we did not really take the
momentum into account. In fact, the truck can accelerate because it pushes
back the whole earth behind it!
So, let us take up the computation with a terrestrial mass M , which is
large but not infinite. We will take the limit [M ] at the very end of the
computation, and more precisely, we will let [M /m ].
At the beginning of the experiment, in the terrestrial reference frame,
the speed of the truck is v. At the end of the experiment, the speed is v .
Earth, on the other hand, has original speed V = 0, and if one remains
m
in the same Galilean reference frame, final speed V = M
( v v ) (because
of conservation of total momentum).1 The kinetic energy of the system at

the beginning is then 12 m v2 and at the end it is 12 m v 2 + 12 M V 2 . So, the
difference is given by
h
m i
1
1 m2
1
E c = m( v 2 v 2 ) +
( v v )2 = m( v 2 v2 ) 1 + O
.
2
2M
2
M
This is the amount of gas involved! So we see that, up to negligible terms, the
first argument gives the right answer, namely, 0.14 gallons.
We now come back to the patrolmans frame, moving with speed w with
respect to the terrestrial frame. The initial speed of the truck is v = v w,
and the final speed is v = v w. The Earth has initial speed V = w
m
and final speed V = w + M
( v v ). The difference is now:
1

1
E c = m v2 v 2 + M V 2 V 2
2
2
i2 1
1
1 hm
1
= m( v w)2 m( v w)2 + M
( v v ) w M w2
2
2
2
M
2
1
1
1 m2
2
2
2
= m v m v + m( v v ) w +
( v v ) m( v v ) w,
2
2
2M
E c = E c .
Hence the difference of kinetic energy is preserved, as we expected. So even

in this other reference frame, a correct computation shows that the quantity
of gas involved is the same as before.
The patrolmans mistake was to forget the positive term m( v v ) w,
corresponding to the difference of kinetic energy of the Earth in its Galilean
frame. This term does not tend to 0 as [M /m ] !
From the point of view of the patrolmans frame, 0.02 gallons are needed to
accelerate the truck, and the remaining 0.12 gallons are needed to accelerate the Earth!
We can summarize this as a table, where T is the truck and E is the Earth.
Initial
speed
Final
speed
Ec init.
Ec final
Ec
1
m v2
2
1
m v2
2
m 2
( v v2 )
2
m2
+
( v v )2
2M
v w
m
M
( v v )
v w
m
M
( v v ) w
1 m2
(v
2 M
0
1
m( v
2
1
2
1
m( v
2
w)2
M w2
v )2
M
2
m
M
w)2
( v v ) w
2
m 2
( v v2 )
2
m2
+
( v v )2
2M
1
Note that the terrestrial reference frame is then not Galilean, since the Earth started
moving under the trucks impulsion.
(Illustration
c Craig Thompson 2001)
Fig. 1.1 Romeo, Juliet, and the boat on a lake.
The second paradox is resolved in the same manner: the Earths rebound
energy must be taken into account after the shock with the ball.
The interested reader will find another paradox, relating to optics, in Exercise 1.3 on page 43.
1.1.b
Romeo, Juliet, and viscous fluids
Here is an example in mechanics where a function f (x) is defined on [0, +[,

but lim x0+ f (x) 6= f (0).
Let us think of a summer afternoon, which Romeo and Juliet have dedicated to a pleasant boat outing on a beautiful lake. They are sitting on each
side of their small boat, immobile over the waters. Since the atmosphere is
conducive to charming murmurs, Romeo decides to go sit by Juliet.
Denote by M the mass of the boat and Juliet together, m that of Romeo,
and L the length of the walk from one side of the boat to the other (see
Figures 1.1 and 1.2).
Two cases may be considered: one where the friction of the boat on the
lake is negligible (a perfect fluid), and one where it is given by the formula
f = v, where f is the force exerted by the lake on the boat, v the speed of
the boat on the water, and a viscosity coefficient. We consider the problem
only on the horizontal axis, so it is one-dimensional.
We want to compute how far the boat moves
1. in the case = 0;
2. in the case 6= 0.
Let be this distance.
The first case is very easy. Since no force is exerted in the horizontal plane
on the system boat + Romeo + Juliet, the center of gravity of this system
(Illustration
c Craig Thompson 2001)
Fig. 1.2 Romeo moved closer.
does not move during the experiment. Since Romeo travels the distance L
relative to the boat, it is easy to deduce that the boat must cover, in the
opposite direction, the distance
m
=
L.
m+M
In the second case, let x(t) denote the position of the boat and y(t) that
of Romeo, relative to the Earth, not to the boat. The equation of movement
for the center of gravity of the system is
M x + m y =
x.
We now integrate on both sides between t = 0 (before Romeo starts moving)
and t = +. Because of the friction, we know that as [t +], the speed
of the boat goes to 0 (hence also the speed of Romeo, since he will have been
long immobile with respect to the boat). Hence we have
+

(M x + m y )
= 0 = x(+) x(0)
0
or = 0. Since 6= 0, we have = 0, whichever way Romeo moved to the

other side. In particular, if we take the limit when 0, hoping to obtain
the nonviscous case, we have:
lim () = 0
0
>0
hence
lim () 6= (0).
0
>0
V (x)
V0
x
Fig. 1.3 Potential wall V (x) = V0 H (x).
Conclusion: The limit of a viscous fluid to a perfect fluid is singular. It is not

possible to formally take the limit when the viscosity tends to zero to obtain
the situation for a perfect fluid. In particular, it is easier to model flows of
nonviscous perfect fluids by real fluids which have large viscosity, because of
turbulence phenomena which are more likely to intervene in fluids with small
viscosity. The interested reader can look, for instance, at the book by Guyon,
Hulin, and Petit [44].
Remark 1.1 The exact form f = v of the friction term is crucial in this argument. If the
force involves additional (nonlinear) terms, the result is completely different. Hence, if you try
to perform this experiment in practice, it will probably not be conclusive, and the boat is not
likely to come back to the same exact spot at the end.
1.1.c Potential wall in quantum mechanics

In this next physical example, there will again be a situation where we have a
limit lim f (x) 6= f (0); however, the singularity arises here in fact because of
x0
a second variable, and the true problem is that we have a double limit which
does not commute: lim lim f (x, y) 6= lim lim f (x, y).
x0 y0
y0 x0
The problem considered is that of a quantum particle arriving at a potential wall. We look at a one-dimensional setting, with a potential of the type
V (x) = V0 H (x), where H is the Heaviside function, that is, H (x) = 0 if
x < 0 and H (x) = 1 if x > 0. The graph of this potential is represented in
Figure 1.3.
A particle arrives from x = in an energy state E > V0 ; part of it is
transmitted beyond the potential wall, and part of it is reflected back. We are
interested in the reflection coefficient of this wave.
The incoming wave may be expressed, for negative values of x, as the sum
of a progressive wave moving in the direction of increasing x and a reflected
wave. For positive values of x, we have a transmitted wave in the direction
of increasing x, but no component in the other direction. According to the
Schrdinger equation, the wave function can therefore be written in the form
(x, t) = (x) f (t),
V (x)
V0
Fig. 1.4 Smoothed potential V (x) = V0 /(1 + e x/a ), with a > 0.
where
e ik x + B e ik x
(x) =
ik x
Ae
def
if x < 0, with k =
if x > 0, with
def
p
2mE
,
h}
2m(E V0 )
.
h}
The function f (t) is only a time-related phase factor and plays no role in
what follows. The reflection coefficient of the wave is given by the ratio of
the currents associated with and is given by R = 1 kk |A|2 (see [20, 58]).
Now what is the value of A? To find it, it suffices to write the equation
expressing the continuity of and at x = 0. Since (0+ ) = (0 ), we
have 1 + B = A. And since (0+ ) = (0 ), we have k(1 B) = k A, and
we deduce that A = 2k/(k + k ). The reflection coefficient is therefore equal
to
2
p
p

E E V0
k
k k
2
R = 1 |A| =
= p
.
(1.1)
p
k
k + k
E + E V0
Here comes the surprise: this expression (1.1) is independent of h}. In particular, the limit as [ h} 0] (which defines the classical limit) yields a nonzero
reflection coefficient, although we know that in classical mechanics a particle
with energy E does not reflect against a potential wall with value V0 < E !2
So, displaying explicitly the dependency of R on h}, we have:
lim R( h}) 6= 0 = R(0).
h}0
h}6=0
In fact, we have gone a bit too fast. We take into account the physical
aspects of this story: the classical limit is certainly not the same as brutally
writing h} 0. Since Plancks constant is, as the name indicates, just a
constant, this makes no sense. To take the limit h} 0 means that one
arranges for the quantum dimensions of the system to be much smaller than
all other dimensions. Here the quantum dimension is determined by the de
2
The particle goes through the obstacle with probability 1.
Broglie wavelength of the particle, that is, = h}/ p. What are the other
lengths in this problem? Well, there are none! At least, the way we phrased it:
because in fact, expressing the potential by means of the Heaviside function is
rather cavalier. In reality, the potential must be continuous. We can replace it
by an infinitely differentiable potential such as V (x) = V0 /(1 + e x/a ), which
increases, roughly speaking, on an interval of size a > 0 (see Figure 1.4). In
the limit where a 0, the discontinuous Heaviside potential reappears.
Computing the reflection coefficient with this potential is done similarly,
but of course the computations are more involved. We refer the reader to [58,
chapter 25]. At the end of the day, the reflection coefficient is found to
depend not only on h}, but also on a, and is given by

sinh a(k k ) 2
R( h}, a) =
.
sinh a(k + k )
p
p
( h} appears in the definition of k = 2mE / h} and k = 2m(E V0 )/ h}.)
We then see clearly that for fixed nonzero a, the de Broglie wavelength of the
particle may become infinitely small compared to a, and this defines the
correct classical limit. Mathematically, we have
a 6= 0
lim R( h}, a) = 0
classical limit
h}0
h}6=0
On the other hand, if we keep h} fixed and let a to to 0, we are converging

to the Heaviside potential and we find that

k k 2
h} 6= 0
lim R( h}, a) =
= R( h}, 0).
a0
k + k
a6=0
So the two limits [ h} 0] and [a 0] cannot be taken in an arbitrary order:

k k 2
lim lim R( h}, a) =
but
lim lim R( h}, a) = 0.
a0 h}0
h}0 a0
k + k
h}6=0 a6=0
a6=0 h}6=0
To speak of R(0, 0) has a priori no physical sense.
1.1.d Semi-infinite filter behaving as waveguide

We consider the circuit A B,
2C
A
L
made up of a cascading sequence of T cells (2C , L, 2C ),
10

2C
2C
A
L
B
(the capacitors of two successive cells in series are equivalent with one capacitor
with capacitance C ). We want to know the total impedance of this circuit.
First, consider instead a circuit made with a finite sequence of n elementary
cells, and let Zn denote its impedance. Kirchhoffs laws imply the following
recurrence relation between the values of Zn :
1
iL
+
Z
n
1
2iC
+
,
(1.2)
Zn+1 =
1
2iC
+ Zn
iL +
2iC
where is the angular frequency. In particular, note that if Zn is purely
imaginary, then so is Zn+1 . Since Z1 is purely imaginary, it follows that
Zn iR
for all n N.
We dont know if the sequence (Zn )nN converges. But one thing is certain: if this sequence (Zn )nN converges to some limit, this must be purely
imaginary (the only possible real limit is zero).
Now, we compute the impedance of the infinite circuit, noting that this
circuit A B is strictly equivalent to the following:
2C
2C
A
L
Hence we obtain an equation involving Z:
1
iL
+
Z
1
2iC
Z=
+
.
1
2iC
+Z
iL +
2iC
(1.3)
Some computations yield a second-degree equation, with solutions given by
1
1
Z2 = L
.
C
4C 2
We must therefore distinguish two cases:
11
1
If < c = p
, we have Z 2 < 0 and hence Z is purely imaginary of
2 LC
the form
r
1
L
Z = i
.
2
2
4C
C
Remark 1.2 Mathematically, there is nothing more that can be said, and in particular there
remains an uncertainty concerning the sign of Z .

However, this can also be determined by a physical argument: let tend to 0 (continuous
regime). Then we have
i
,
Z () +
0
2C
the modulus of which tends to infinity. This was to be expected: the equivalent circuit is
open, and the first capacitor cuts the circuit. Physically, it is then natural to expect that, the
first coil acting as a plain wire, the first capacitor will be dominant. Then
Z () +
0
i
2C
(corresponding to the behavior of a single capacitor).

Thus the physically acceptable solution of the equation (1.3) is
r
L
1
.
Z = i
4C 2 2
C
1
If > c = p
, then Z 2 > 0 and Z is therefore real:
2 LC
r
L
1
Z =
.
C
4C 2 2
Remark 1.3 Here also the sign of Z can be determined by physical arguments. The real part
of an impedance (the resistive part) is always non-negative in the case of a passive component,
since it accounts for the dissipation of energy by the Joule effect. Only active components
(such as operational amplifiers) can have negative resistance. Thus, the physically acceptable
solution of equation (1.3) is
r
L
1
.
Z= +
C
4C 2 2
In this last case, there seems to be a paradox since Z cannot be the limit
as n + of (Zn )nN . Lets look at this more closely.
From the mathematical point of view, Equation (1.3) expresses nothing but the
fact that Z is a fixed point for the induction relation (1.2). In other words,
this is the equation we would have obtained from (1.2), by continuity, if we
had known that the sequence (Zn )nN converges to a limit Z. However, there
is no reason for the sequence (Zn )nN to converge.
Remark 1.4 From the physical point of view, the behavior of this infinite chain is rather sur-
prising. How does resistive behavior arise from purely inductive or capacitative components?
Where does energy dissipate? And where does it go?
12
In fact, there is no dissipation of energy in the sense of the Joule effect, but energy does
disappear from the point of view of an operator holding points A and B. More precisely,
one can show that there is a flow of energy propagating from cell to cell. So at the beginning
of the circuit, it looks like there is an energy sink with fixed power consumption. Still, no
energy disappears: an infinite chain can consume energy without accumulating it anywhere.3
In the regime considered, this infinite chain corresponds to a waveguide.
We conclude this first section with a list of other physical situations where
the problem of noncommuting limits arises:
taking the classical (nonquantum) limit, as we have seen, is by no
means a trivial matter; in addition, it may be in conflict with a nonrelativistic limit (see, e.g., [6]), or with a low temperature limit;
in plasma thermodynamics, the limit of infinite volume (V ) and
the nonrelativistic limit (c ) are incompatible with the thermodynamic limit, since a characteristic time of return to equilibrium is
V 1/3 /c ;
in the classical theory of the electron, it is often said that such a classical
electron, with zero naked mass, rotates too fast at the equator (200 times
the speed of light) for its magnetic moment and renormalized mass to
conform to experimental data. A more careful calculation by Lorentz
himself4 gave about 10 times the speed of light at the equator. But in
fact, the limit [m 0+ ] requires care, and if done correctly, it imposes
a limit [v/c 1 ] to maintain a constant renormalized mass [7];
another interesting example is an infinite universe limit and a diluted

universe limit [56].
1.2
Sequences
1.2.a
Sequences in a normed vector space
We consider in this section a normed vector space (E , kk) and sequences of

elements of E .5
DEFINITION 1.5 (Convergence of a sequence) Let (E , kk) be a normed vec-
tor space and (un )nN a sequence of elements of E , and let E . The
3
This is the principle of Hilberts infinite hotel.

Pointed out by Sin-Itiro Tomonaga [91].
5
We recall the basic definitions concerning normed vector spaces in Appendix A.
4
Sequences
13
sequence (un )nN converges to if, for any > 0, there exists an index
starting from which un is at most at distance from :
> 0
N N
n N
n N = kun k < .
Then is called the limit of the sequence ( u n ) nN

N , and this is denoted
= lim un
n
or
un .
n
DEFINITION 1.6 A sequence (un )nN of real numbers converges to + (resp.
to ) if, for any M R, there exists an index N , starting from which all
elements of the sequence are larger than M :
M R
N N
n N
n N = un > M
(resp. un < M ).
In the case of a complex-valued sequence, a type of convergence to infinity,

in modulus, still exists:
DEFINITION 1.7 A sequence (zn )nN of complex numbers converges to infinity if, for any M R, there exists an index N , starting from which all
elements of the sequence have modulus larger than M :
M R
N N
n N
n N = |zn | > M .
Remark 1.8 There is only one direction to infinity in C. We will see a geometric interpretation of this fact in Section 5.4 on page 146.
Remark 1.9 The strict inequalities kun k < (or |un | > M ) in the definitions above (which,
in more abstract language, amount to an emphasis on open subsets) may be replaced by

kun k , which are sometimes easier to handle. Because > 0 is arbitrary, this gives
an equivalent definition of convergence.
1.2.b Cauchy sequences

It is often important to show that a sequence converges, without explicitly
knowing the limit. Since the definition of convergence depends on the limit ,
it is not conveninent for this purpose.6 In this case, the most common tool
is the Cauchy criterion, which depends on the convergence of elements of the
sequence with respect to each other:
DEFINITION 1.10 (Cauchy criterion) A sequence (un )nN in a normed vector
space is a Cauchy sequence, or satisfies the Cauchy criterion, if

> 0 N N p, q N
q > p N = u p uq <
6
As an example: how should one prove that the sequence (un )nN , with
n
P
(1) p+1
un =
,
p4
p=1
converges? Probably not by guessing that the limit is 74 /720.
14
or, equivalently, if
> 0 N N p, k N

p N = u p+k u p < .
A common technique used to prove that a sequence (un )nN is a Cauchy

sequence is therefore to find a sequence ( p ) pN of real numbers such that

lim p = 0
and
p, k N u p+k u p p .
p
PROPOSITION 1.11 Any convergent sequence is a Cauchy sequence.
This is a trivial consequence of the definitions. But we are of course

interested in the converse. Starting from the Cauchy criterion, we want to be
able to conclude that a sequence converges without, in particular, requiring
the limit to be known beforehand. However, that is not always possible:
there exist normed vector spaces E and Cauchy sequences in E which do not
converge.
Example 1.12 Consider the set of rational numbers Q. With the absolute value, it is a normed
Q-vector space. Consider then the sequence

u0 = 3
u1 = 3.1
u2 = 3.14
u3 = 3.141
u4 = 3.1415
u5 = 3.14159
(you can guess the rest ...). This is a sequence of rationals, which is a Cauchy sequence (the
distance between u p and u p+k is at most 10 p ). However, it does not converge in Q, since its
limit (in R!) is , which is a notoriously irrational number.
The space Q is not nice in the sense that it leaves a lot of room for Cauchy sequences to
exist without converging in Q. The mathematical terminology is that Q is not complete.
DEFINITION 1.13 (Complete vector space) A normed vector space (E , kk) is

complete if all Cauchy sequences in E are convergent.
THEOREM 1.14 The spaces R and C, and more generally all finite-dimensional real
or complete normed vector spaces, are complete.
Proof.
First case: It is first very simple to show that a Cauchy sequence (un )nN of real numbers is bounded. Hence, according to the Bolzano-Weierstrass theorem (Theorem A.41,
page 581), it has a convergent subsequence. But any Cauchy sequence which has a
convergent subsequence is itself convergent (its limit being that of the subsequence),
see Exercise 1.6 on page 43. Hence any Cauchy sequence in R is convergent.
Second case: Considering C as a normed real vector space of dimension 2, we can
suppose that the base field is R.
Consider a basis B = (b 1 , . . . , b d ) of the vector space E . Then we deduce that E is
complete from the case of the real numbers and the following two facts: (1) a sequence
( x n )nN of vectors, with coordinates (xn1 , . . . , xnd ) in B, converges in E if and only if
each coordinate sequence (x k n )nN ; and (2), if a sequence is a Cauchy sequence, then
each coordinate is a Cauchy sequence.
7
This is simply un = 10n 10n where is the integral part function.
Sequences
15
Both facts can be checked immediately when the norm of E is defined by k xk =

max|x k |, and other norms reduce to this case since all norms are equivalent on E .
Example 1.15 The space L 2 of square integrable functions (in the sense of Lebesgue), with the
def R
norm k f k2L2 = R | f |2 , is complete (see Chapter 9). This infinite-dimensional space is used
very frequently in quantum mechanics.
Counterexample 1.16 Let E = K[X ] be the space of polynomials

P with coefficients in K (and
arbitrary degree). Let P E be a polynomial, written as P = n X n , and define its norm by
def
kP k = max |i |. Then the normed vector space (E , kk) is not complete (see Exercise 1.7 on
page 43).iN
Here is an important example of the use of the Cauchy criterion: the fixed
point theorem.
1.2.c The fixed point theorem

We are looking for solutions to an equation of the type
f (x) = x,
where f : E E is an application defined on a normed vector space E , with
values in E . Any element of E that satisfies this equation is called a fixed
point of f .
DEFINITION 1.17 (Contraction) Let E be a normed vector space, U a subset
of E . A map f : U E is a contraction if there exists a real number

[0, 1[ such that f ( y) f (x) k y xk for all x, y U . In particular,
f is continuous on U .
THEOREM 1.18 (Banach fixed point theorem) Let E be a complete normed vec-
tor space, U a non-empty closed subset of E , and f : U U a contraction. Then f

has a unique fixed point.
Proof. Chose an arbitrary u0 U , and define a sequence (un )nN for n 0 by

induction by un+1 = f (un ). Using the definition of contraction, an easy induction
shows that we have

u p+1 u p p ku1 u0 k
for any p 0. Then a second induction on k 0 shows that for all p, k N, we have
p

u p+k u p ( p + + p+k1 ) ku1 u0 k ku1 u0 k ,
1
and this proves that the sequence (un )nN is a Cauchy sequence. Since the space E is
complete, this sequences has a limit a E . Since U is closed, we have a U .
Now from the continuity of f and the relation un+1 = f (un ), we deduce that
a = f (a). So this a is a fixed
point of f . If b is an arbitrary fixed point, the inequality
ka b k = f (a) f (b ) ka b k proves that ka b k = 0 and thus a = b ,
showing that the fixed point is unique.
16
Remark 1.19 Here is one reason why Banachs theorem is very important. Suppose we have
a normed vector space E and a map g : E E , and we would like to solve an equation
g(x) = b . This amounts to finding the fixed points of f (x) = g(x) + x b , and we can hope
that f may be a contraction, at least locally. This happens, for instance, in the case of the
Newton method, if the function used is nice enough, and if a suitable (rough) approximation
of a zero is known.
This is an extremely fruitful idea: one can prove this way the Cauchy-Lipschitz theorem
concerning existence and unicity of solutions to a large class of differential equations; one can
also study the existence of certain fractal sets (the von Koch snowflake, for instance), certain
stability problems in dynamical systems, etc.
Not only does it follow from Banachs theorem that certain equations have
solutions (and even better, unique solutions!), but the proof provides an effective way to find this solution by a successive approximations: it suffices to fix
u0 arbitrarily, and to define the sequence (un )nN by means of the recurrence
formula un+1 = f (un ); then we know that this sequence converges to the fixed
point a of f . Moreover, the convergence of the sequence of approximations
is exponentially fast: the distance from the approximate solution un to the
(unknown) solution a decays as fast as n . An example (the Picard iteration) is
given in detail in Problem 1 on page 46.
1.2.d Double sequences

Let (xn,k )(n,k)N2 be a double-indexed sequence of elements in a normed vector
space E . We assume that the sequences made up of each row and each column
converge, with limits as follows:
x11
x21
x31
..
.
B1
x12
x22
x32
..
.
B2
x13 A1
x23 A2
x33 A3
..
.
B3
The question is now whether the sequences (An )nN and (Bk )kN themselves
converge, and if that is the case, whether their limits are equal. In general, it
turns out that the answer is No. However, under certain conditions, if one
sequence (say (An )) converges, then so does the other, and the limits are the
same.
DEFINITION 1.20 A double sequence (xn,k )n,k converges uniformly with respect to k to a sequence (B k ) kN
N as n if

> 0 N N n N k N
n N = xn,k Bk < .
In other words, there is convergence with respect to n for fixed k, but in such
a way that the speed of convergence is independent of k; or one might say that
all values of k are similarly behaved.
Sequences
17
Uniform convergence with respect to n toward the sequence (An )nN is

similarly defined.
THEOREM 1.21 (Double limit) With notation as above, if the following three con-
ditions hold:
each row converges, and An = lim xn,k for all n N,

k
each column converges, and Bk = lim xn,k for all k N,

n
the convergence is uniform either with respect to n or with respect to k;

then the sequences (An )nN and (Bk )kN converge, lim An = lim Bk = . One
n
says that the double sequence ( x n, k ) n, k converges to the limit .
Be aware that the uniform convergence condition is very important.

The following examples gives an illustration: here both limits exist, but they
are different.
1 0 0 0 0
1 1 0 0 0
1 1 1 0 0
1 1 1 1 0
.. .. .. ..
..
.
. . . .

1 1 1 1
1.2.e
Sequential definition of the limit of a function
DEFINITION 1.22 Let f : K K (where K, K = R or C or any normed
vector space), let a K, and let K . Then f has the limit , or tends
to , at the point a if we have

> 0 > 0 z K
|z a| < = f (z) < .
There are also limits at infinity and infinite limits, defined similarly:
DEFINITION 1.23 Let f : K K (where K, K = R or C). Let K .
Then f tends to at +, resp. at (in the case K = R), resp. at infinity

(in the case K = C), if we have

> 0 A R x R
x > A = f (x) < ,

resp. > 0 A R x R
x < A = f (x) <

resp. > 0 A R+ z C
|z| > A = f (z) < .
Similarly, a function f : R R tends to + at + if
M > 0 A R
x R
x > A = f (x) > M ,
18
and finally a function f : C C tends to infinity at infinity if

M > 0 A R z C
|z| > A = f (z) > M .
In some cases, the definition of limit is refined by introducing a punctured

neighborhood, i.e., looking at the values at points other than the point where
the limit is considered:
DEFINITION 1.24 Let f : K {a} K (with K, K = R or C and a K)
and let K . Then f ( x) converges to in punctured neighborhoods

of a if
> 0
z K

(z 6= a and |z a| < ) = f (z) < .
This is denoted
= za
lim f (z).
z6=a
This definition has the advantage of being practically identical to the definition of convergence at infinity. It is often better adapted to the physical
description of a problem, as seen in Examples 1.1.b and 1.1.c on page 5 and
the following pages. A complication is that it reduces the applicability of the
theorem of composition of limits.
THEOREM 1.25 (Sequential characterization of limits) Let f : K K be a
function, and let a K and K . Then f (x) converges to as x tends to a if

and only
if, for any convergent sequence (xn )nN with limit equal to a, the sequence
f (xn ) nN converges to .
1.2.f
Sequences of functions
Consider now the case of a sequence of functions ( fn )nN each defined on a

same subset X of R or C and taking values in R or C. Denote by kk the
supremum norm8:

k f k = sup f (x).
xX
The definition of convergence of real or complex sequences may be extended

to functions in two ways: one is a local notion, called simple, or pointwise,
convergence, and the other, more global, is uniform convergence.9
8
In fact, it is not a norm on the space of all functions, but only on the subspace of
bounded functions. We disregard this subtlety and consider here that kk takes values in R+ =
R {+}.
9
The concept of uniform convergence is due to George Stokes (see page 472) and Philipp
Seidel (18211896), independently.
Sequences
19
DEFINITION 1.26 (Simple convergence) Let ( fn )nN be a sequence of functions, all defined on the same set X , which may be arbitrary. Then the
sequence ( f n ) nN
N converges simply (or: pointwise on
X) to a function f
defined on X if, for any x in X , the sequence fn (x) nN converges to f (x).
This is denoted
cv.s.
fn f .
DEFINITION 1.27 (Uniform convergence) Let ( fn )nN be a sequence of func-
tions, all defined on the same set X , which may be arbitrary. Then the
sequence ( f n ) nN
N converges uniformly to the function f if
> 0 N N
This is denoted
n N
n N = k fn f k < .
cv.u.
fn f .
In other words, in the case where X is a subset of R, the graph of the

function f is located inside a smaller and smaller band of constant width in
which all the graphs of fn must also be contained if n is large enough:
If we have functions fn : X E , where (E , kk) is a normed vector space,

we similarly define pointwise and uniform convergence using convergence
in E ; for instance, ( fn )nN converges uniformly to f if

> 0 N N n N
n N = sup fn (x) f (x) < .
xX
Remark 1.28 Uniform convergence is an important theoretical mathematical notion. If one

wishes to compute numerically a function f using successive approximations fn (for instance,
partial sums of a series expansion), then to get an error of size at most for the
value f (x1 )
of f at some given point x1 , it suffices to find N1 such that fn (x1 ) f (x1 ) for any
n N1 . Similarly, if the value of f at other points x2 , x3 , . . . , x p is needed, it will be enough
to find corresponding integers N2 , N3 , . . . , N p . However, if it is not known beforehand at
which points the function will evaluated, it will be necessary to know an integer N such that
fn (x) f (x) for all n N and for all x R. Uniform convergence is then desirable.
It is clear that uniform convergence implies pointwise convergence, but the

converse is not true.
Example 1.29 Define a sequence ( fn )n1 a functions on R by
nx
fn (x) = 2 nx
if x [0, 1/n] ,
if x [1/n, 2/n] ,
if x [2/n, 1] .
1/n 2/n
20
The reader will have no trouble proving that ( fn )n1 converges pointwise to the zero function.
However, the convergence is not uniform, since we have k fn f k = 1 for all n 1.
Example 1.30 The sequence of functions fn : R R defined for n 1 by

x
fn : x 7 sin x +
n
converges uniformly to f : x 7 sin x on the interval [0, 2], and in particular it converges
pointwise on this interval. However, although the sequence converges pointwise to the sine
function on R, the convergence is not uniform on all of R. Indeed, for n 1, we have
n
n
n
n
fn
= sin
+
and
f
= sin
,
2
2
2
2
2
and those two values differ by 1 in absolute value. However, one can check that the convergence
is uniform on any bounded segment in R.
2
Exercise 1.1 Let g(x) = e x and fn (x) = g(x n). Does the sequence ( fn )nN converge
pointwise on R? Does it converge uniformly?
Remark 1.31 In the case where the functions fn are defined on a subset of R with finite
measure (for instance, a finite segment), a theorem of Egorov shows that pointwise convergence
implies uniform convergence except on a set of arbitrarily small measure (for the definitions,
see Chapter 2).
Remark 1.32 There are other ways of defining the convergence of a sequence of functions. In
particular, when some norm is defined on a function space containing the functions fn , it is
possible do discuss convergence in the sense of this norm. Uniform convergence corresponds
to the case of the kk norm. In Chapter 9, we will also discuss the notion of convergence in
quadratic mean, or convergence in L2 norm, and convergence in mean or convergence in L1 norm.
In pre-Hilbert spaces, there also exists a weak convergence, or convergence in the sense of scalar
product (which is not defined by a norm if the space is infinite-dimensional).
A major weakness of pointwise convergence is that it does not preserve

continuity (see Exercise 1.10 on page 44), or limits in general (Exercise 1.12).
Uniform convergence, on the other hand, does preserve those notions.
THEOREM 1.33 (Continuity of a limit) Let ( fn )nN be a sequence of functions de-
fined on a subset D in K (or in a normed vector space), with values in an arbitrary

normed vector space. Assume that the sequence ( fn )nN converges uniformly to a function f .
i) Let a D be such that all functions fn are continuous at the point a. Then f
is also continuous at the point a.
ii) In particular, if each fn is continuous on D , then the limit function f is also
continuous on D .
This property extends to the case where a is not in D , but is a limit point
of D . However, it is then necessary to reinforce the hypothesis to assume that
the functions fn have values in a complete normed vector space.
Sequences
21
THEOREM 1.34 (Double limit) Let D be a subset of R (or of a normed vector

space) and let x0 D be a limit point10 of D . Let ( fn )nN be a sequence of functions
defined on D with values in a complete normed vector space E . Assume that, for all
n, the function fn has a limit as x tends to x0 . Denote n = lim xx0 fn (x).
If ( fn )nN converges uniformly to a function f , then
i) f (x) has a limit as x x0 ;
ii) (n )nN has a limit as n ;

iii) the two limits are equal: lim f (x) = lim n .
xx0
In other words:
lim lim f (x)

xx0 n n
= lim lim fn (x).

n xx0
If we want a limit of differentiable functions to remain differentiable,

stronger assumptions are needed:
THEOREM 1.35 (Differentiation of a sequence of functions) Let I be an interval of R with non-empty interior, and let ( fn )nN be a sequence of functions defined
on I with values in R, C, or a normed vector space. Assume that the functions fn are
differentiable on I , and moreover that:
i) the sequence ( fn )nN converges pointwise to a function f ;

ii) the sequence ( fn )nN converges uniformly to a function g.
Then f is differentiable on I and f = g.
Remark 1.36 If the functions take values in R or C or more generally any complete normed
vector space, it is possible to weaken the assumptions by asking, instead of (i), that the sequence
( fn (x0 )) converges at a single point x0 I . Assumption (ii) remains identical, and the conclusion is the same: ( fn )nN converges uniformly to a differentiable function with derivative equal
to g.
Counterexample 1.37 The sequence of functions given by
fn : x 7
n
8 X
sin2 (nx)
k=1 4n2 1
converges uniformly to the function f : x 7 |sin x| (see Exercise 9.3 on page 270), but the
sequence of derivatives does not converge uniformly. The previous theorem does not apply,
and indeed, although each fn is differentiable at 0, the limit f is not.
Remark 1.38 It happens naturally in some physical situations that a limit of a sequence of
functions is not differentiable. In particular, in statistical thermodynamics, the state functions
of a finite system are smooth. However, as the number of particles grows to infinity, discontinuities in the state functions or their derivatives may appear, leading to phase transitions.
10
See Definition 4.52 on page 106; the simplest example is D = ]a, b ] and x0 = a.
22
Uniform convergence is also useful in another situation: when trying to

exchange a limit (or a sum) and an integration process. However, in that
situation, pointwise convergence is often sufficient, using the powerful tools
of Lebesgue integration (see Chapter 2).
THEOREM 1.39 (Integration on a finite interval) Let ( fn )nN be a sequence of
integrable functions (for instance, continuous functions), which converges uniformly to

a function f on a finite closed interval [a, b] R. Then f is integrable on [a, b]
and we have
Z
Z
b
lim
fn (x) dx =
f (x) dx.
Example 1.40 This theorem is very useful, for instance, when dealing with a power series
expansion which is known to converge uniformly
open disc of convergence (see TheoP on the
n
rem 1.66 on page 34). So, if we have f (x) =
n=0 a n x for |x| < R, then for any x such that
|x| < R, we deduce that
Zx
X
an
x n+1 .
f (s) ds =
n
+1
0
n=0
To establish that a sequence converges uniformly, in practice, it is necessary

to compute k fn f k , or rather to bound this expression by a quantity which
itself converges to 0. This is sometimes quite tricky, and it is therefore useful
to know the following two results of Dini:11
THEOREM 1.41 (Dini) Let K be a compact subset of Rk , for instance, a closed ball.
Let ( fn )nN be an increasing sequence of continuous functions converging pointwise

on K to a continuous function f . Then the sequence ( fn )nN converges uniformly to f
on K .
THEOREM 1.42 (Dini) Let I = [a, b] be a compact interval in R, and let ( fn )nN
be a sequence of increasing functions from I to R that converges pointwise on I to a

continous function f . Then ( fn )nN converges uniformly on I .
Remark 1.43 Be careful to distinguish between an increasing sequence of functions and
a sequence of increasing functions. The former is a sequence ( fn )nN of real-valued

functions such that fn+1 (x) fn (x) for any x R and n N. The latter is a
sequence of real-valued functions defined on K R such that for any x, y K :
x y = fn (x) fn ( y).
As mentioned briefly already, it is possible with Lebesgues dominated

convergence theorem to avoid requiring uniform convergence to exchange an
integral and a limit, as in Theorem 1.39. See Chapters 2 and 3 for details on
this theory.
11
Ulisse Dini (18451918) studied in Pisa and Paris before taking a position in Pisa. His
work concerned the theory of functions of a real variable, and he contributed to the early
development of functional analysis.
Series
23
1.3
Series
1.3.a
Series in a normed vector space
We first recall the definition of convergence and absolute convergence of a

series in a normed vector space.
DEFINITION 1.44 (Convergence of a series) Let (an )n be a sequence with val-
ues in a normed vector space. Let (Sn )nN denote the sequence of partial
sums
n
def P
ak .
Sn =
k=0
The series
a n converges, and its sum is equal to A if the sequence
(Sn )nN converges to A. This is denoted
n=0
The series
in R.
an = A.
a n converges absolutely if the series
kan k converges
P
In particular, a series
P an of real or complex numbers converges absolutely if the series |an | converges in R.
As in the case of sequences, there exists a Cauchy criterion for convergence

of series12:
P
THEOREM 1.45 (Cauchy criterion) If a series
un with values in a normed vector space E converges, then it satisfies the Cauchy criterion:

X

q

un
> 0 N N p, q N
(q > p N ) =
< ,
n= p
or in other words:
lim
p,q
un = 0.
pnq
Conversely, any series which satisfies the Cauchy criterion and takes values in R,
C, any finite-dimensional normed vector space, or more generally, any complete normed
vector space, converges.
12
This criterion was stated by Bernhard Bolzano (see page 581) in 1817. But Bolzano was
isolated in Prague and little read. Cauchy presented this criterion, without proof and as an
obvious fact, in his analysis course in 1821.
24
From this the following fundamental theorem is easily deduced:

P
THEOREM 1.46 Any absolute convergent series
an with values in a complete
normed vector space is convergent.
In particular, any absolutely convergent series of real or complex numbers is convergent.
Proof. Let
un be an absolutely convergent series. Although we can write

X
X
k
k

u
kun k ,
n

n=0
n=0
nothing can be deduced from this, because the right-hand side does not tend to zero.
But we can use the Cauchy critetion: for all p, q N, we have of course
X
X
q
q

u
kun k ,
n

n= p
n= p
P
and since
kun k satisfies the Cauchy criterion, so does
P un . Since un lies in a
complete space by assumption, this means that the series un is indeed convergent.
1.3.b Doubly infinite series
In the theory of Fourier series, we will have to deal with formulas of the type
Z1
+
X

f (t)2 dt =
|cn |2 .
0
n=
To give a precise meaning to the right-hand side, we must clarify the meaning
of the convergence of a series indexed by integers in Z instead of N.
P
DEFINITION 1.47 A doubly
series nZZ a n , with an is a normed
P infinite P
vector space, converges if
an and
an are both convergent, the index
ranging over N in each case. Then we denote
+
X
n=
and say that this is the sum of
def
an =
P
an +
n=0
nZ a n .
X
n=1
an
P
In other words, a series of complex numbers nZ an converges to if and
only if, for any > 0, there exists N > 0 such that
X

j

for any i N and j N ,
an < .

n=i
Remark 1.48 It is crucial to allow the

Pk upper and lower bounds i and j to be indepen-
dent. In particular,
P if the limit of
infinite series nZ an converges.
n=k
an exists, it does not follow that the doubly
Series
25
P
For instance, take an = 1/n for n 6= 0 and a0 = 0. Then we have kn=kP
an = 0
for all k (and so this sequence does converge as k tends to infinity),
but
the
series
nZ a n
P
P
diverges according to the definition, because each of the series an and an (over n 0) is
divergent.
1.3.c
Convergence of a double series
As in Section 1.2.d, let (a i j ) i, jN be a family of real numbers indexed by two

integers. For any p, q N, we have
p X
q
X
ai j =
i=1 j=1
q X
p
X
ai j ,
j=1 i=1
since each sum is finite. On the other hand, even if all series involved are
convergent, it is not always the case that
X
X
ai j
X
X
and
i=1 j=1
ai j ,
j=1 i=1
as the following example shows:
1 1
0
0
0 1 1
0
(a i j ) = 0 0
1 1
..
..
..
..
.
.
.
.
0 ...
0 . . .
0 . . .
..
.
0
0
0
..
.
where we have (note that i is the row index and j is the column index):
X
ai j = 0
but
i=1 j=1
We can find an even more striking example by
1 1
0
0
0
0 2 2
0
0
(a i j ) = 0 0
3 3
0
..
..
..
.. ..
.
.
.
.
.
in which case
X
X
i=1 j=1
ai j =
X
i=1
a i j = 1.
j=1 i=1
0=0
but
putting
X
X
j=1 i=1
0 ...
0 . . .
0 . . .
..
.
ai j =
X
j=1
1 = +.
26
1.3.d Conditionally convergent series,

absolutely convergent series
P
an with an in a normed vector space is conditionnally convergent if it is convergent but not absolutely convergent.
DEFINITION 1.49 A series
DEFINITION 1.50 We denote by S the group of permutations, that is, the

group of bijections from N to N, and we denote by Sn the finite group of
permutations of the set {1, . . . , n}.
P
DEFINITION 1.51 A series
x n is commutatively convergent if it is convergent
with
sum
X
,
and
for
any
permutation S, the rearranged series
P
x(n) converges to X .
Is a convergent series necessarily commutatively convergent? In other words,

is it legitimate to change arbitrarily the order of the terms of a convergent series?
At first sight, it is very tempting to say Yes, almost without thinking,
since permuting terms in a finite sum has no effect on the result. The problem
is that we have here an infinite sum, not a finite sum in the algebraic sense.
So there is a limiting process involved, and we will see that this brings a
very different picture: only absolutely convergent series will be commutatively
convergent. So, if a series is conditionally convergent (convergent but not
absolutely so), changing the order of the terms may alter the value of the sum
or even turn it into a divergent series.
P
THEOREM 1.52 Let
an be a conditionnally convergent series, with an R. Then
for
P any R, there exists a permutation S such that the rearranged series
a(n) converges to .
In fact, for any a, b R, with a b, there exists a permutation
Pn S such
that
the set of limit points of the sequence of rearranged partial sums
k=0 a (k) nN is
the interval [a, b].
P
Proof. We assume that an is conditionnally convergent.
First remark: there are infinitely many positive values and infinitely many negative values
of the terms an of the series. Let n denote the sequence of non-negative terms, in the
order they occur, and let n denote the sequence of negative terms.
Here is an illustration:
un : 1
3
2 4 1 2 1 0
2
1 2 3 1 2 4 3 5 6
P
P
Second remark: both series n and
n are divergent, their partial sums converging,
P respectively, to + and . Indeed, if both series were to converge, the
series
an would be absolutely convergent, and if only one were to converge, then
P
an would be divergent (as follows from considering the sequence of partial sums).
Third remark: both sequences (n )nN andP(n )nN tend to 0 (since (an )nN tends
to 0, as a consequence of the convergence of an ).
Now consider R. Let Sn denote the sequence of sums of values of and
which is constructed as follows. First, sum all consecutive values of n until their sum
Series
27
is larger than ; call this sum S1 . Now add to S1 all consecutive values of n until
the resulting sum S1 + 1 + is smaller than ; call this sum S2 . Then start again
adding from the remaining values of n until getting a value larger than , called S3 ,
and continue in this manner until the end of time.
Now notice that:
Since at each step we add at least one value of or one of , it is clear that
all values of will be used sooner or later, as well as all values of , that is,
when all is said and done, all values of an will have been involved in one of the
sums Sn .
Since, at each step, the distance | Sn | is at most equal to the absolute value
of the last value of or considered, the distance from Sn to tends to 0 as n
tends to infinity.
From this we deduce that
P the sequence (Sn ) is a sequence of partial sums of a
rearrangement of the series an , and that it converges to . Hence this proves that by
simply changing the order of the terms, one may cause the series to converge to an arbitrary sum.
Let now a, b R with a < b (the case a = b being the one already considered).
If a and b are both finite, we can play the same game of summation as before,
but this time, at each step, we either sum values of n until we reach a value
larger than b , or we sum values of n until the value is less than a.
If b = + and a is finite, we sum from a to above a + 1, then come back to

below a, then sum until we are above a + 2, come back below a, etc. Similarly
if a = and b is finite.
If a = and b = +, start from 0 to go above 1, then go down until

reaching below 2, then go back up until reaching above 3, etc.
Example 1.53 Consider the sequence (an )nN with general term an = (1)n+1
P/n. It follows
from the theory of power series (Taylor expansion of log(1 + x)) that the series an converges
and has sum equal to log 2. If we sum the same values an by taking one positive term followed
by two negative terms, then the resulting series converges to 12 log 2. Indeed, if (Sn )nN and
(S n )nN denote the sequence of partial sums of the original and modified series, respectively,
then for n N we have
1 1 1
1
1
S2n = 1 + + +
2 3 4
2n 1 2n
and
1 1 1 1 1
1
1
1
S3n
= + + +
2
4
3
6
8
2n
1
4n
2
4n
|{z}
| {z }
|
{z
}
=
1 1 1 1
1
1
1
+ + +
= S2n .
2 4 6 8
4n 2 4n
2
As an exercise, the reader can check that if one takes instead two positive terms followed by
one negative terms, the resulting series converges with a value equal to 32 log 2.
The following result shows that, on the other hand, one can rearrange at
will the order of the terms of an absolutely convergent series.
THEOREM 1.54 A series of complex numbers is commutatively convergent if and only
if it is absolutely convergent.
28

Proof. Assume
P first that the terms of the series are real numbers. The theorem above
shows that if
an is commutatively convergent, it must be absolutely convergent.
Conversely, assume the series is absolutely convergent and let denote its sum. Let
be any permutation. By convergence of the series, there exists N 1 such that
X

n

<
a
k

k=1
for n N . For each such n N , there exists N such that the set {(1), . . . , (N )}
contains {1, . . . , n} (it suffices that N be larger than the maximum of the images of 1,
. . . , N by the inverse permutation 1 ). Then for any m N , we have
X
X
X
X
m
n

+
a
|ak | +
|ak |,
(k)
k

k=1
k=1
k>n
k>n
since the set {(1), . . . , (m)} contains {1, . . . , n}, and possibly additional values which
are all larger than n. The absolute convergence makes its P
appearance now: the last sum
on the right is the remainder for the convergent series |ak |, and for n N it is
therefore itself smaller than . Since, given , we can take n = N and find the value
N from it, such that
X

m

2
a
(k)

k=1
for m N , and so we have proved that the rearranged series converges with sum equal
to .
If the terms of the series are complex numbers, it suffices to apply the result for real
series to the series of real and imaginary parts.
The possibility of rearranging at will the order of summation explains the

importance of absolutely convergent series13.
Remark 1.55 In statistical mechanics, there are so-called diagrammatic methods to compute the
values of certains quantities, such as pressure or mean-energy, at equilibrium. Those methods

are based on rearrangements of the terms of certain series, summing by packets in particular.
Those methods are particularly useful when the original series is not absolutely convergent. This
means that all results obtained in this manner must be treated carefully, if not suspiciously.
They belong to the gray area of exact (at least, this is what everyone believes!) results, but which
are not rigorous. (It is of course much more difficult to obtain results which can be judged with
mathematical standards of rigor; the reader is invited to read the beautiful papers [63, 64] for
convincing illustrations.)
13
Peter Gustav Lejeune-Dirichlet showed in 1837 that a convergent series with non-negative
terms is commutatively convergent. In 1854, Bernhard Riemann wrote three papers in order to
obtain a position at the university of Gttingen. In one of them, he describes commutatively
convergent series in the general case. However, another paper was selected, concerning the
foundations of geometry.
Series
29
1.3.e Series of functions

We can define pointwise and uniform convergence of series of functions just
as was done for sequences of functions.
DEFINITION 1.56 (Pointwise convergence)
Let X be an arbitrary set, (E , kk)
P
a normed vector space. A series

fn of functions fn : X E converges
P
pointwise to a function F : X E if, for any x X , the series
fn (x)
converges to F (x) in E , that is, if
> 0 x X
N N
n N
n

P

n N =
fk (x) F (x)
< .
k=1
P
The function F is called the pointwise, or simple, limit of the series
fn ,
P cv.s.
and this is denoted
fn F .
DEFINITION 1.57 (Uniform convergence)

P
a normed vector space. A series
uniformly to a function F : X E if the sequence of partial sums of the
series converges uniformly to F , that is, if
> 0 N N
x X
n N

n

P
fk (x) F (x)
n N =
< .

k=1
cv.u.
This is denoted
fn F . This amounts to

n

P

lim
fk F
=0
where
kgk = sup g(x) .

n
k=1
xX
DEFINITION 1.58 (Absolute convergence)

P
a normed vector space.PA series
absolutely if the series k fn k converges, where

k fn k = sup fn (x) .
xX
The following theorem is the most commonly used to prove uniform

convergence of a series of functions:
THEOREM 1.59 Any absolutely convergent series with values in a complete normed
vector space is uniformly convergent, and hence pointwise convergent.
Corresponding to the continuity and differentiability results for sequences of

functions, we have:
30
THEOREM 1.60 (Continuity and differentiability of a series of functions)
Let D be a subset of R or of a normed

vector space. Let ( fn )nN be a sequence of
functions fn : D E , where
E
,
kk
is some normed vector space, for instance, R
P
or C. Assume that the series
fn converges pointwise to a function F .
P
i) If each fn is continuous on D , and if the series
fn converges uniformly on
D , then F is continuous on D .
P
ii) If D is an interval of R, each fn is differentiable on D , and the series
fn
converges uniformly, then F is differentiable and we have
F =
fn .
n=0
1.4
Power series, analytic functions
Quite often, physicists encounter series expansions of some function.
These expansions may have different origins:
the superposition of many phenomena (as in the Fabry-Perot interferometer);
perturbative expansions, when exact computations are too difficult to
perform (e.g., hydrodynamics, semiclassical expansions, weakly relativistic expansions, series in astronomy, quantum electrodynamics, etc.);
sometimes the exact evalution of a function which expresses some physical quantity is impossible; a numerical evaluation may then be performed using Taylor series expansions, Fourier series, infinite product
expansions, or asymptotic expansions.
We first recall various forms of the Taylor formula. The general idea is
that there is an approximate expression
(x a)2
(x a)k (k)
f (a) + +
f (a)
2!
k!
for a function f which is at least k times differentiable
on an interval J , with

values in some normed vector space E , kk , and for a given point a J ,
where x lies is some neighborhood of a.
The question is to make precise the meaning of the symbol !
Define Rk (x) to be the difference between f (x) and the sum on the righthand side of the above expression; in other words, we have
f (x) f (a) + (x a) f (a) +
f (x) =
k
X
(x a)n
n=0
n!
f (n) (a) + Rk (x) = Tk (x) + Rk (x)
31
Brook Taylor (16851731), English mathematician, was a student

at Cambridge, then member and secretary of the prestigious Royal
Society, a venerable institution dedicated to the advancement of
Science. He wrote the famous formula
f (a + ) = f (a) + f (a) +
2
2
f (a) +
without considering the issue of convergence. He was also interested in the physical and mathematical aspects of vibrating strings.
by definition. Of course, we hope that the Taylor remainder R k ( x) is a

small quantity, so that we may approximate the value of f (x) by the value
of the Taylor polynomial of order k at x, that is, by Tk (x). There are
different ways in which this remainder may become small:
one may let x tend to a (for a fixed value of k);
or let k tend to infinity (for a fixed value of x).
The Taylor-Lagrange and Taylor-Young formulas are relevant for the first case,
while the second belongs to the theory of power series.
1.4.a
Taylor formulas
THEOREM 1.61 (Taylor

formula with integral remainder) Let J be an interval

of R, and E , kk a normed vector space. Let f : J E be a function of C k class
on J , which is piecewise of C k+1 class on J . For any a and x J , we have
Zx
k
X
(x a)n (n)
(x t)k (k+1)
f (x) =
f (a) +
f
(t) dt.
n!
k!
a
n=0
THEOREM 1.62 (Taylor-Lagrange formula) Let f : J R be a real-valued
function of C k class on an interval J of R, which is k + 1 times differentiable in the

interior of J . Let a J . Then, for any x J , there exists ]0, 1[ such that
f (x) =
k
X
(x a)n
n=0
n!
f (n) (a) +

(x a)k+1 (k+1)
f
a + (x a) .
(k + 1)!
Remark 1.63 This formula is only valid for real-valued functions. However, the following
corollary is also true for functions with complex values, or functions with values in a normed
vector space.
32
COROLLARY 1.64 (Taylor-Lagrange inequality) Let f : J E be a function of
C k class on an interval J of R, with values in a normed vector space E . Assume f

is k + 1 times differentiable in the interior of J . Let a J . Then for any x J we
have

k
k+1
X

(k+1)
(x a)n (n)
f (x)
|x a|
f
f
(a)
sup
(t) E .

n!
(k
+
1)!
t
J
E
n=0
THEOREM 1.65 (Taylor-Young formula) Let f be a function which is k times
differentiable on an interval J of R, with values in a normed vector space E . Let
a J . Then we have
f (x)
k
X
(x a)n
n!
n=0
1.4.b

f (n) (a) = o (x a)k .
xa
Some numerical illustrations
Suppose we want to compute numerically some values of the inverse tangent

function arctan, which is of course infinitely differentiable on R. It is easy to
compute the values of the successive derivatives of this function at 0, and we
can write down explicitly the Taylor polynomial at 0 of arbitrary order: this
gives the expression
arctan x =
k
X
(1)n 2n+1
x
+ Rk (x),
2n + 1
n=0
for the Taylor formula of order 2n + 1 (notice that only odd powers of x
appear, because the inverse tangent function is odd).
If we represent graphically those polynomials with k = 0, 1, 4, 36 (i.e., of
order 1, 5, 9, and 18, respectively), with the graph of the function itself for
comparison, we obtain the following:
0.5
0.5
0.5
0.5
0.5
0.5
0.5
0.5
0.5
0.5
0.5
0.5
0.5
0.5
0.5
0.5
order 9
order 37
order 1
order 5
1
x
The following facts appear:

on each graph (i.e., for fixed k), the Taylor polynomial and the inverse
tangent functions get closer and closer together as x approaches 0;
33
for a fixed real number x [1, 1] (for instance 0.8), the values at x of
the Taylor polynomial of increasing degree get closer and closer to the
value of the function as k increases;
on the other hand, for a real number x such that |x| > 1, disaster
strikes: the larger k is, the further away to arctan x is the value of the
Taylor polynomial!
The first observation is simply a consequence of the Taylor-Young formula.
The other two deserve more attention. It seems that the sequence (Tk )kN of
the Taylor polynomials converges on [1, 1] and diverges outside.14 However,
the function arctan is perfectly well-defined, and very regular, at the point x = 1; it
does not seem that anything special should happen there. In fact, it is possible
to write down the Taylor expansion centered at a = 1 instead of a = 0 (this
is a somewhat tedious computation15), and (using approximations of the same
order as before), we obtain the following graphs:
2
We can see the same three basicp

facts, except
p that convergence seems to be
restricted now to the interval 1 2, 1 + 2 .
In order to understand why such intervals occur, it is necessary to dwell further on the theory of power series (see below) and especially on holomorphic
functions of a complex variable (in particular, Theorem 4.40 on page 101).
We will only state here that the function arctan can be continued naturally
to a function on the complex plane (arctan z0 is defined as the value of the
integral of the function 1/(1 + z 2 ) on a certain path16 joining the origin to z0 .
The function thus obtained is well-defined, independently of the chosen path,
up to an integral multiple of 17 and is a well-defined function on C minus
the two single points where 1 + z 2 vanishes, namely i and i. Then, one
shows that for such a function, the sequence of Taylor polynomials centered
14
To be honest, it is difficult to ascertain from the graphs above if the interval to consider is
[1, 1] or ]1, 1[, for instance. The general theory of series shows that (Tn (x))nN converges
quickly if |x| < 1 and very slowly if |x| = 1.
15
The n-th coefficient of the polynomial is (1)n+1 sin(n/4)2n/2 /n and the constant term
is /4.
16
This notion of integral on a path is defined by the formula (4.2) page 94.
17
This is a consequence of the residue theorem 4.81 on page 115.
34
at a converges on the open disc centered at a with radius equal to the distance
from a to the closest singulirity (hence the radius isp|1 0| = 1 in the first
case of Taylor expansions at a = 0, and is |1 i| = 2 in the second case).
1.4.c Radius of convergence of a power series

A power series centered at z0 is any series of functions of the type
z 7
X
n=0
an (z z0 )n ,
where (an )nN is a given sequence of real or complex numbers, which are
sometimes called the coefficients of the power series.
P
THEOREM 1.66 (Radius of convergence) Let
an (z z0 )n be a power series
+
centered at z0 . The radius of convergence is the element in R defined by

def
R = sup t R+ ; (an t n )nN is bounded .
The power series converges absolutely and uniformly on any compact subset in the disc

def
B(z0 ; R) = z C ; |z z0 | < R , in the complex plane C, and it diverges for
any z C such that |z| > r. For |z| = r, the series may be convergent, conditionally
convergent, or divergent at z.
Note that absolute convergence here refers to absolute convergence as a
series of functions, which is stronger than absolute convergence for every z
involved: in other words, for any compact set D B(z0 ; R), we have
sup |an z n | < +.
n=0 zD
Example 1.67 The power series log(1 z) =
n
n=1 z /n converges for any z C such that
|z| < 1 and diverges if |z| > 1 (the radius of convergence is R = 1). Moreover, this series is
divergent at z = 1, but conditionnally convergent at z = 1 (by the alternate series test), and
more generally, it is conditionnally convergent at z = e i for any
/ 2Z (this can be shown
using the Abel transformation, also known as summation by parts).
DEFINITION 1.68 (Power series expansion) Let be an open subset in R

or C. A function f : C defined on has a power series expansion
centered at some z0 if there exist an open subset V containing z0
and a sequence (an )nN of complex numbers such that
z V
f (z) =
X
n=0
an (z z0 )n .
The radius of convergence of a power series depends only weakly on the

precise values of the coefficients, so, for instance, if F = P / Q is a rational
35
P
P
function with no pole in N, the power series
F (n) an z n and an z n have
the same radius of convergence. From this and Theorem 1.60, it follows in
particular that a power series can be differentiated term by term inside the
disc of convergence:
THEOREM 1.69 (Derivative of a power series) Let J be an open subset of R,
x0 J , and f : J C a function which has a power series expansion centered at x0 :

f (x) =
n=0
an (x x0 )n .
Let R > 0 be the radius of convergence of this power series. Then f is infinitely differentiable on the open interval ]x0 R, x0 + R[, and each derivative has a power series
expansion on this interval, which is obtained by repeated term by term differentiation,
that is, we have
f (x) =
X
n=1
nan (x x0 )n1
f (k) (x) =
and
X
n=k
n!
a (x x0 )nk
(n k)! n
for any k N. Hence the n-th coefficient of the power series f (x) can be expressed as
an =
f (n) (x0 )
.
n!

f (n) (x0 )/n! (x x0 )n is the Taylor series of f at x0 . On
any compact subset inside the open interval of convergence, it is the uniform limit of the
sequence of Taylor polynomials.
Remark 1.70 The power series
Remark 1.71 In Chapter 4, this result will be extended to power series of one complex variable
(Theorem 4.40 on page 101).
1.4.d
Analytic functions
Consider a function that may be expended into a power series in a neighborhood V of a point z0 , so that for z V , we have
P
f (z) =
an (z z0 )n .
n=0
Given such a z V , a natural question is the following: may f also be expended

into a power series centered at z?
Indeed, it might seem possible a priori that f can be expanded in power
series only around z0 , and around no other point. However, this is not the
case:
DEFINITION 1.72 (Analytic function) A function f : C defined on an
open subset of C or R is analytic on if, for any z0 , f has a power
series expansion centered at z0 .
36
Note that the radius of convergence of the power series may (and often
does!) vary with the point z0 .
P
THEOREM 1.73 Let
an z n be a power series with positive radius of convergence
R > 0, and let f denote the sum of this power series on B(0 ; R). Then the function
f is analytic on B(0 ; R).
Example 1.74 The function f : x 7 1/(1 x) has the power series expansion
f (x) =
xn
n=0
around 0, with radius of convergence equal to 1. Hence, for any x0 ]1, 1[, there exists a
power series expansion centered at x0 (obviously with different coefficients). This can be made
explicit: let h B(0 ; |1 x0 |), then with x = x0 + h, we have
f (x) = f (x0 + h) =
X
1
1
1
(x x0 )n
=
=
.
1 (x0 + h)
1 x0 1 h/(1 x0 )
(1 x0 )n+1
n=0
Remark 1.75 (Convergence of Taylor expansions) Let f : U C be a function defined on an
open subset U of R. Under what conditions is f analytic?18 There are two obvious necessary
conditions:
f is infinitely differentiable on U ;
for any x0 U , there exists an open disc B(x0 , r ) such that the series
x0 )n converges for any x B(x0 , r ).
1
n!
f (n) (x0 )(x
However, those two conditions are not sufficient. The following classical counter-example
shows this: let

def
f (x) = exp 1/x 2
if x 6= 0,
f (0) = 0.
It may be shown19 that f is indeed of C and that each derivative of f at 0 vanishes, which
ensures (!) the convergence of the Taylor series everywhere. But since the function vanishes
only at x = 0, it is clear that the Taylor series does not converge to f on any open subset,
hence f is not analytic.
It is therefore important not to use the terminology analytic where infinitely differentiable is intended. This is a confusion that it still quite frequent in scientific literature.
The Taylor formulas may be used to prove that a function is analytic. If the sequence
(Rn )nN of remainders for a function f converges uniformly to 0 on a neighborhood of a R,
then the function f is analytic on this neighborhood. To show this, one may use the integral
expression of the remainder terms in the Taylor formula. A slightly different but useful
approach is to prove that both the function under consideration and its Taylor series (which
must be shown to have positive radius of convergence) satisfy the same differential equation,
with the corresponding initial conditions; then f is analytic because of the unicity of solutions
to a Cauchy problem.
Also, it is useful to remember that if f and g have power series expansions centered at z0 ,
then so do f + g and f g. And if f (z0 ) 6= 0, the function 1/ f also has a power series expansion
centered at z0 .
18
The same question, for a function of a complex variable, turns out to have a completely
different, and much simpler, answer: if f is differentiable in the complex sense on the
open set of definition, then it is always analytic. See Chapter 4.
19
By induction, proving that f (n) (x) is for x 6= 0 of the form x 7 Q n (x) f (x), for some
rational function Q n .
A quick look at asymptotic and divergent series
37
1.5
1.5.a
Asymptotic series
DEFINITION 1.76 (Asymptotic expansion) Let F be a function of a real or

complex variable z, defined for all z with |z| large enough. The function F
has an asymptotic expansion if there exists a sequence (an )nN of complex
numbers such that
(
)
N
X
a
n
lim z N F (z)
=0
n
z
z
n=0
for any positive integer N . This is denoted
X
an
F (z)
.
z
zn
n=0
(1.4)
The definition means that the expansion (1.4) is a good approximation for
large values of z. Indeed, if we only consider the first twenty terms of the
series, for instance, we see that the sum of those approximates f (z) to order
1/z 20 at least when [z ].
However, it frequently turns out that for fixed z, the behavior of the series
in (1.4) is quite bad as N . In particular, the series may be divergent.
This phenomenon was pointed out and studied in detail by Henri Poincar
in the case of asymptotic series used in astronomy, at the beginning of the
twentieth century [70].
P
How can a divergent asymptotic series still be used? Since
an /z n is
asymptotic to F , if there is some R such that F is continuous for |z| R,
then we see that there exist constants C1 , C2 ,. . . such that

N
X

a
n
F (z)
CN
for N N and |z| R.

n
z |z| N +1
n=0
For fixed z, we can look for the value of N such that the right-hand side of
this inequality is minimal, and truncate the asymptotic series at this point.
Of course, we do not obtain F (z) with infinite precision. But in many cases the
actual precision increases with |z|, as described in the next section, and may
be pretty good.
It is also possible to speak of asymptotic expansion as z 0, which
corresponds to the existence of a sequence (an )nN such that

N
X
1
n
lim N f (z)
an z = 0,
z0 z
n=0
38
which is denoted
f (z)
z0
an z n .
(1.5)
n=0
Remark 1.77 If it exists, an asymptotic expansion of a function is unique, but there
may be two different functions with the same asymptotic expansion! For instance, e x
2
and e x both have asymptotic expansions with an = 0 for all n as x +.
A physical example is given by quantum electrodynamics. This quantum

theory of electromagnetic interactions gives physical results in the form of
series in powers of the coupling constant = e 2 / h}c 1/137 (the Sommerfeld fine structure constant), which means that a perturbative expansion in is
performed:
As shown by Dyson [32], when studying a physical quantity F we can expect to find a perturbative series of the following type (with the normalization
h} = c = 1):
X
F (e 2 ) = F () =
fn n .
n=0
Since the value of is fixed by Nature, with a value given by experiments, only
the truncated series can give a physical result if the series is divergent. The
truncation must be performed around the 137-th term, which means that we
still expect a very precise result certainly more precise, by far, than anything
the most precise experiment will ever give! However, if F (e 2 ) is not analytic
at e = 0, the question is raised whether the asymptotic expansion considered
gives access to F uniquely or not.
Studying asymptotic series is in itself a difficult task. Their implications
in physics (notably field theory) are at the heart of current research [61].
1.5.b Divergent series and asymptotic expansions

Since Euler, Cauchy (page 88), and especially Poincar (page 475), it has
been realized that divergent series may be very useful in physics. As seen in the
previous section, they appear naturally in computations of asymptotic series.
As a general rule, convergent series are used to prove numerical or functional identities (between power series, Fourier series, etc). Thus the series may
be used instead of the value of their sum at any time in a computation. Another remark is that, from the computational viewpoint, some series are more
interesting than others, because they converge faster. For example, we have the
following two identities for log 2:
1 1 1
(1)n+1
+ + +
+
2 3 4
n
1
1
1
1
1
log = log 2 = + 2
+
+ + n
+
2
2 2 2 23 3
2 n
The second of those (which comes from expanding x 7 log(1 x) in power
series at x = 1/2) converges much faster than the first (which results from a
log 2 = 1
39
Leonhard Euler (17071783), a Swiss mathematician, an exceptional teacher, obtained a position at the Academy of Sciences of
Saint Petersburg thanks to Nicolas and Daniel Bernoulli when
he was only twenty. He also spent some years in Berlin, but
came back to Russia toward the end of his life, and died there at
seventy-six (while drinking tea). His works are uncountable! We
owe him the notations e and i and he imposed the use of that
was introduced by Jones in 1706. Other notations due to Euler
are sin, cos, tang, cot, sec, and cosec. He also introduced the use
of complex exponents, showed that e i x = cos x + i sin x, and was
particularly fond of the formula e i + 1 = 0. He defined the
function , which extends the factorial function from integers
to C \ (N), and used the Riemann zeta function for real values
of the variable. No stone of the mathematical garden of his time
was left unturned by Euler; let us only add the Euler angles in
mechanics and the Euler equation in fluid mechanics.
similar expansion at x = 1, hence on the boundary of the disc of convergence20).

While studying problems of celestial mechanics, Poincar realized that the
meaning of convergent series was not the same for mathematicians, with
rigor in mind, or astronomers, interested in efficiency:
Geometers, preoccupied with rigorousness and often indifferent to the
length of the inextricable computations that they conceive, with no idea of
implementing them in practice, say that a series is convergent when the sum
of its terms tends to some well-defined limit, however slowly the first terms
might diminish. Astronomers, on the contrary, are used to saying that a series converges when the twenty first terms, for instance, diminish very quickly,
even though the next terms may well increase indefinitely. Thus, to take a
1000n
simple example, consider the two series with general terms
and
1 2 3n
1 2 3n
.
1000n
Geometers will say that the first series converges, and even that it converges
20
The number of terms necessary to approximate log 2 within 106 , for instance, can be
estimated quite precisely for both series. Using the Leibniz test for alternating sums, the
remainder of the first series is seen to satisfy

|Rn | un+1 =
1
,
n+1
and this is the right order of magnitude (a pretty good estimate is in fact Rn 1/2n). If we
want |Rn | to be less than 106 , it suffices to take n = 106 terms. This is a very slow convergence.
The remainder of the second series, on the other hand, can be estimated by the remainder of
a geometric series:
X
X
1
1
1
Rn =
= n.
n
k
n
2
2
2
k=n+1
k=n+1
Hence twenty terms or so are enough to approximate log 2 within 106 using this expansion
(since 220 106 ).
40
Fig. 1.5 The precise value of f (x)Pis always found between two successive values of the
partial sums of the serie
fk (x). Hence, it is inside the gray strip.
rapidly, [...] but they will see the second as divergent.
Astronomers, on the contrary, will see the first as divergent [...] and the
second as convergent. [70]
How is it possible to speak of convergence of a series which is really

divergent? We look at this question using a famous example, the Euler series.
Let
Z + t/x
e
dt,
x > 0
f (x) =
1+t
0
and say we wish to study P
the behavior of f for small values of x. A first idea
is to expand 1/(1 + t) as (1)k t k and exchange the sum and the integral if
permitted. Substitute y = t/x and then integrate by parts; a simple induction
then shows that
Z +
t k e t/x dt = k! x k+1 ,
0
P
and since the (1)k k! x k+1 is obviously divergent for any nonzero value
of x, this first idea is a lamentable failure.
To avoid this problem, it is possible to truncate the power-series expansion
of the denominator, and write
n1
X
(1)n t n
1
=
(1)k t k +
,
1+ t
1+t
k=0
from which we derive an expression for f of the type f = fn + Rn , where

fn (x) = x x 2 + 2! x 3 3! x 4 + + (1)n1 (n 1)! x n
(1.6)
Z + n t/x
and
t e
Rn (x) = (1)n
dt.
1+ t
0

Since (1 + t)1 1, the remainder satisfies Rn (x) n! x n+1 , which means
that Rn (x) is of absolute value smaller than the first omitted term; moreover,
41
0,056
0,054
0,052
0,05
0,048
0,046
0,044
10
20
30
40
50
P
Fig. 1.6 The first 50 partial sums of the series (1)k1 (k 1)! x k for x = 1/20.
Notice that, starting from k = 44, the series diverges rapidly. The best precision
is obtained for n = 20, and gives f (x) with an error roughly of size 2 108 .
they are of the same sign. It follows (see the proof of the alternating series
test) that
(1.7)
f2n (x) < f (x) < f2n+1 (x),
although, in contrast with the case of alternating series with terms converging
to 0, the general term here (1)n n! x n+1 diverges. Hence it is not possible to
deduce from (1.7) an arbitrarily precise approximation of f (x). However, if x is
small, we can still get a very good approximation, as we now explain.
of x. There exists an index N0 such that the distance
Fix a positive value

f2n+1 (x) f2n (x) is smallest (the ratio between consecutive terms is equal to
nx, so this value of n is in fact N0 = 1/x). This means that, if we look
at the first N0 values, the series seems to converge, before it starts blowing
up. It is interesting to remark that the convergence
of the first

N0 terms
is exponentially fast, since the minimal distance f N +1 (x) f N (x) is roughly
given by
p
N ! x N N ! N N 2/x e 1/x
(using the Stirling formula, see Exercise 5.4 on page 154.) Thus, if we wish
to know the
p value of f (x) for a small value of x, and if a precision of the
order of 2/x e 1/x suffices, it is possible to use the divergent asymptotic
series (1.6), by computing and summing the terms up to the smallest term (see
Figure 1.5). For instance, we obtain for x = 1/50 a precision roughly equal
to 6 1020 , which is perfectly sufficient for most physical applications! (see
Figure 1.6.)
For a given value of x, on the other hand, the asymptotic series does not
allow any improvement on the precision.21 But the convergence is so fast that
21
For instance, in quantum field theory, the asymptotic series in terms of has a limited
42

Sir George Biddel Airy (18011892), English astronomer, is known
in particular for discovering the theory of diffraction rings. He
determined approximately the solar apex, the direction toward
which the sun and the solar system seem to be directed, in the
Hercules region. He was also interested in geology. He was
director of the Royal Observatory and took part in the controversy concerning priority for the discovery of Neptune (the
French pushing the claim of Le Verrier while the English defended Adams).
The picture here represents a fake stamp, painted directly on the
envelope, representing a contemporary caricature of Sir Airy; the
post office was bluffed and stamped and delivered the letter.
it makes it possible to do some computations which are out of reach of a

standard method! And what Poincar remarked is, in fact, a fairly general
rule: divergent series converge, in general, much more rapidly than convergent series.
In 1857, George Stokes was studying the Airy integral [3]
Z +
3

t
def 1
Ai(z) =
cos
+ z t dt,
0
3
which appears in the computations of caustics. The goal was to find zeros of
this function, and compare the with experimental zeros (corresponding to
dark bands in any optics figure, which had been measured with great precision, at least as far as the first twenty-five). Airy himself, using a convergent
series expansion of Ai at 0, managed fairly easily to compute the position of
the first band, and with considerable difficulty, found the second one. In
fact, his mathematically convergent expansion was divergent in the sense of
astronomers (all the more so as one gets farther away from the origin). Stokes
used instead the devilish method of divergent series22 and, after bypassing
some nontrivial difficulties (linked, in particular, to complex integration), obtained all the hands23 with a precision of 104 !
Remark 1.78 There are other well-known techniques to give a sense to the sum of (some)
divergent series. The interested reader may read the classic book of mile Borel [13], the first
part of which at least is very readable. Concerning asymptotic expansions, see [72].
precision since is fixed (equal to 1/137 approximately) and cannot be made to tend to zero.
This suggests that quantum field theory, in its current perturbative form, will one day be
replaced by another theory. Of course, as long as a precision to 10100 is enough...
22
Niels Abel wrote in 1826 that divergent series are the Devils invention, and it is shameful
to base any proof of any kind on them. By using them, one can get from them whatever
result is sought: they have done much evil and caused many paradoxes (letter to his professor
Holmbo).
23
Only the first is less precise, because it is too small and Stokes used an asymptotic expansion at +.
Exercises
43
EXERCISES
Physical paradoxes
Exercise 1.2 (Electrical energy) Consider an electric circuit consisting of two identical capac-
itors in series, with capacitance C and resistance R. Suppose that for t 0, the circuit is open,
one of the capacitors carries the charge Q , and the other has no charge. At t = 0, the circuit
is closed, and is left to evolve freely. What is the state of equilibrium for this circuit? What
is the energy of the system at t = 0? What is the energy as t +? Show that the missing
energy depends only on R. What happened to this energy?
Now assume that R = 0. What is the energy of the system at any arbitrary t? What is the
limit of this energy as t +? Do you have any comments?
Exercise 1.3 (A paradox in optics) We know that two distinct sources of monochromatic
light do not create a clear interference picture in an experiment with Young slits. As the
distance between the sources increases, we first see a contrast decrease in the interference picture. This is called a defect of spatial coherence.
Hence, a famous experiment gives a measurement of the angular distance between two
components of a double star by the observation of the first disappearance of the interference
fringes when slowly moving two Young slits apart.
This experiment works very well with monochromatic light. However, if we define two
monochromatic sources S1 and S2 mathematically, each emits a signal proportional to e 2i t ,
and there should be no problem of spatial coherence.
Perform the computation properly. A computation in optics always starts with amplitudes
(possibly, one may show that the crossed terms cancel out in average, and do the computations
with intensity only). Here, the cross terms are fine, and never disappear. In other words, this
shows that two different monochromatic light sources are always perfectly coherent.
But experiment shows the opposite: a defect of spatial coherence. How can this be explained?
Exercise 1.4 In the rubber ball paradox of page 2, give an interpretation of the variation of
kinetic energy of the ball, in the moving reference frame, in terms of the work of the force
during the rebound. The shock may be modeled by a very large force lasting a very short
amount of time, or one can use the formalism of distributions (see Chapter 7).
Sequences and series

Exercise 1.5 It is known that Q, and hence also Q [0, 1], is countable. Let (xn )nN be a
sequence of rational numbers such that Q [0, 1] = {xn ; n N}. Show that the sequence
(xn )nN diverges.
Exercise 1.6 In an arbitrary normed vector space, show that a Cauchy sequence which has
a convergent subsequence is convergent.

Exercise 1.7 Show that the space K[X ] of polynomials with coefficients in K is not complete with the norm given by
P =
n
P
i=1
def
i X i 7 kP k = max |i | .
1in
Exercise 1.8 (Fixed point) Let a, b R be real numbers with a < b , and let f : [a, b ]
[a, b ] be a continuous function with a fixed point . Assume that there exists a real number
44
and an interval V = [ , + ] around , contained in [a, b ] and stable under f (i.e.,

f (x) V if x V ), such that

f (x) f () |x |2 .
x V
i) Let V be such that ( ) < 1. Let u be the sequence defined by induction by
u0 = ,
Show that (un )nN converges to .

ii) Show in addition that |un | 1 ( )
Exercise 1.9 Let
for all n N.
un+1 = f (un )
2n
3
if 0 x 1/2n,
2n x
2
fn (x) = n 2n3 (x 1/2n) if 1/2n x 1/n,
0
if 1/n x 1,
for n N and x [0, 1]. Plot a graph of f , and compute

Z1
Z1
lim
fn (x) dx
and
lim fn (x) dx.
n
Exercise 1.10 Let ( fn )nN be a sequence of functions converging simply to a function f . If
each fn is increasing, show that f is also increasing. Show that the same stability holds for the
properties fn is convex and fn is k-Lipschitz. Show that, on the other hand, it is possible
that each fn is continuous, but f is not (take fn (x) = sin2n x).
Exercise 1.11 Let n be the function defined on [1, 1] by
n (x) =
Z x
0
1 e 1/nt
dt
for n N .
Show that n is infinitely differentiable, and that the sequence (n )nN converges uniformly on [1, 1]. What is its limit?
Let > 0 be given. Show that for any p N, there exists a map p from [1, 1] into R,
infinitely differentiable, such that
( p)
i) (k)
p (0) = 0 for k 6= p, and p (0) = 1.

ii) for k p 1 and x [1, 1], (k)
p (x) .
Now let (an )nN be an arbitrary sequence of real numbers. Construct an infinitely differentiable map f from [1, 1] to R such that f (n) (0) = an for all n N.
P
fn defined on
R, which converges pointwise to a sum F (x), the convergence being uniform on any finite
interval of R, and which, moreover, satisfies:
Exercise 1.12 (Slightly surprising exercise) Construct a series of functions
n N
lim fn (x) = +
x+
but
lim F (x) = .
x+
Exercises
45
Exercise 1.13 Consider a power series centered at the point a C, given by

def
f (z) =
n=0
cn (z a)n .
Let R denote its radius of convergence, and assume R > 0.

i) Prove the Cauchy formula: for any n N and any r ]0, R[, we have
Z 2
1
f (a + r e i ) e i n d.
cn =
2 r n 0
ii) Prove the Gutzmer formula: for r ]0, R[, we have
Z 2
X

1
f (a + r e i )2 d.
|cn |2 r 2n =
2
0
n=0
iii) Prove that if R = +, in which case the sum f (z) of the power series is said to
be an entire function, and if moreover f is bounded on C, then f is constant (this is
Liouvilles theorem, which is due to Cauchy).
iv) Is the sine function a counter-example to the previous result?
Exercise 1.14 Let f be a function of C class defined on an open set R. Show that f
is analytic if and only if, for any x0 , there are a neighborhood V of x0 and positive real
numbers M and t such that
( p)
f (x)
p

x V p N
p! M t .
Function of two variables
Exercise 1.15 Let f : R2 R be a function of two real variables. This exercise gives
examples showing that the limits
lim lim f (x, y)

x0 y0
lim lim f (x, y)

y0 x0
and
lim
(x, y)(0,0)
f (x, y)
are independent: each may exist without the other two existing, and they may exist without
being equal.
xy
if x 2 + y 2 6= 0,
i) Let f (x, y) = x 2 + y 2
0
if x = y = 0.
Show that the limits lim lim f (x, y) and lim lim f (x, y) both exist, but that the
x0 y0
limit24
lim
(x, y)(0,0)
y0 x0
f (x, y) is not defined.
24
The limit of f as the pair (x, y) tends to a value (a, b ) R2 is defined

using any of

the natural norms on R2 , for instance the norm k(x, y)k = max |x| , | y| or the euclidean
p
def
norm k(x, y)k2 = x 2 + y 2 , which are equivalent. Thus, we have
lim
(x, y)(a,b)
(x, y)6=(a,b)
if and only if
> 0 > 0
f (x, y) =

(x, y) 6= (a, b ) and (x a, y b ) < = | f (x, y) | .
46
y + x sin(1/ y) if y 6= 0,
0
if y = 0.
Show that both limits lim f (x, y) and lim lim f (x, y) exist, but on the other
ii) Let f (x, y) =
y0 x0
(x, y)(0,0)
hand lim lim f (x, y) does not exist.

x0 y0

x y + y sin 1
if x =
6 0,
2
2
iii) Let f (x, y) = x + y
x
0
if x = 0.
Show that lim lim f (x, y) exists. Show that neither

x0 y0
x6=0 y6=0
lim
(x, y)(0,0)
f (x, y), nor lim lim f (x, y)

y0 x0
y6=0
exist.
2
2
x y
iv) Let f (x, y) = x 2 + y 2
if x 2 + y 2 6= 0,
if x = y = 0.
Show that the limits lim lim f (x, y) and lim lim f (x, y) both exist, but are different.
x0 y0
y0 x0
PROBLEM
Problem 1 (Solving differential equations) The goal of this problem is to illustrate, in a
special case, the Cauchy-Lipschitz theorem that ensures the existence and unicity of the solution
to a differential equation with a given initial condition.
In this problem, I is an interval [0, a] with a > 0, and we are interested in the nonlinear
differential equation
ty
y =
(E)
1 + y2
with the initial condition
y(0) = 1.
(CI)
The system of two equations (E) + (CI) is called the Cauchy problem. In what follows, E
denotes the space

C (I , R) of real-valued continuous functions defined on I , with the norm
k f k = sup f (t).
tI

i) Let ( fn )nN be a Cauchy sequence in E , kk .

(a) Show that for any x I the sequence fn (x) nN converges in R. For x I , we
def
let f (x) = lim fn (x).

n
(b) Show that ( fn )nN converges uniformly f on I .

(c) Show that the function f : I R is continuous.

(d) Deduce from this that E , kk is a complete normed vector space.
ii) For any f E , define a function ( f ) by the formula

( f ) : I R,
Zt
u f (u)
t 7 ( f )(t) = 1 +
2 du.
0 1 + f (u)
Show that the functions f E which are solutions of the Cauchy problem (E) + (CI)
are exactly the fixed points of .
Solutions of exercises
x
is 1-Lipschitz on R, i.e.,
1 + x2

y
x

1 + y 2 1 + x 2 | y x| .
47
iii) Show that the function x 7
(1.8)
iv) Show that is a contracting map if a is sufficiently small.
v) Show that there exists a unique solution to the Cauchy problem. Give an explicit
iterative method to solve the system numerically (Picard iterations).
Remark 1.79 In general, all this detailed work need not be done: the Cauchy-Lipschitz theorem
states that for any continous function (x, y) which is locally Lipschitz with respect to the
second variable, the Cauchy problem
y = (t, y),
has a unique maximal solution (i.e., a solution defined on a maximal interval).
SOLUTIONS
Solution of exercise 1.2. The energy of the circuit at the beginning of the experiment is
the energy contained in the charged capacitor, namely E = Q 2 /2C . At equilibrium, when
[t ], no current flows, and the charge of each capacitor is Q /2 (it is possible to write
down the necessary differential equations and solve them to check this). Thus the final energy
is E = 2( Q /2)2 /C = E /2. The energy which is dissipated by the Joule effect (computed by
R +
the integral 0 Ri 2 (t) dt, where t 7 i(t) is the current flowing through the circuit at time t)
is of course equal to E E , and does not depend on R.
However, if R = 0, one observes oscillations of charge in each capacitor. The total energy
of the system is conserved (it is not possible to compute it from relations in a quasi-stationary
regime; one must take magnetic fields into account!). In particular, as [t +], the initial
energy is recovered. The explanation for this apparent contradiction is similar to what happened for Romeo and Juliet: the time to reach equilibrium is of order 2/RC and tends to
infinity as [R 0]. This is a typical situation where the limits [R 0] and [t +] do
not commute.
Finally, if we carry the computations even farther, it is possible to take into account the
electromagnetic radiation due to the variations of the electric and magnetic fields. There is
again some loss of energy, and for [t +], the final energy E E = E /2 is recovered.
Solution of exercise 1.3. Light sources are never purely monochromatic; otherwise there
would indeed be no spatial coherence problem. What happens is that light is emitted in wave
packets, and the spectrum of the source necessarily has a certain width > 0 (in a typical
example, this is order of magnitude = 1014 s1 , corresponding to a coherence length of a
few microns for a standard light-bulb; the coherence length of a small He-Ne laser is around
thirty centimeters, and that of a monomode laser can be several miles). All computations must
be done first with 6= 0 before taking a limit 0. Thus, surprisingly, spatial coherence
is also a matter of temporal coherence. This is often hidden, with the motto being since the
sources are not coherent, I must work by summing intensities instead of amplitudes.
In fact, when considering an interference figure, one must always sum amplitudes, and then
(this may be a memory from your optics course, or an occasion to read Born and Wolf [14])
perform a time average over a period t, which may be very small, but not too much (depending
on the receptor; the eyes are pretty bad in this respect, an electronic receptor is better, but none
can have t = 0).
48
The delicate issue is to be careful of a product t . If you come to believe (wrongly!)

that two purely monochromatic sources interfere without any spatial coherence defect, this
means that you have assumed t = 0. To see the spatial coherence issue arise, one must
keep large enough so that t cannot be neglected in the computation.
P
Solution of exercise 1.7. Let Pn = nk=1 X k /k. It is easy to see that (Pn )nN is a Cauchy
1
sequence: for any integers k and p, we have P p+k P p = p+1
.
However, this sequence does not converge in K[X ]; indeed, if it were to converge to a
limit L, we would have L K[X ], and all coefficients of L of degree large enough ( N ,
say) would be zero, which implies that kPk Lk 1/(deg L + 1) if k N , contradicting that
(Pn )nN converges to L.
This example shows that K[X ] is not complete.
R
Solution of exercise
1.9. For any n N, we have fn = n/2, and for any x [0, 1], the

sequence fn (x) n tends to 0, showing that
Z
Z

lim fn (x) dx = +
whereas
lim fn (x) dx = 0.
n
Solution of exercise 1.11. The sequence (n )nN converges uniformly to 0.

Notice also the property
n (0) = 1,
n (0) = 0,
n(k) (0) = 0
k 2.
For given > 0, it suffices to define p as the (p 1)-st primitive of N , where N is

sufficiently large so that sup (pp1) (x) . Here, each primitive is selected to be the
x[1,1]
one vanishing at 0 (i.e., the integral from 0 to x of the previous one). It is easy to see that
the successive derivatives of this function satisfy the require condition, and the last property
follows from the construction.
Now let (an )nN be an arbitrary sequence of real numbers. For all n N, one can apply the
previous construction to find n such that

1
.
sup (n1)
(x) n
n
2 max(1, |an |)
x[1,1]
P
It is then immediate that the series
an n converges uniformly to a function f having all
desired properties.
Of course, the function f thus constructed is by no means unique: one may add a term
(n 1), where R, without changing the values of the derivatives at 0.
Solution of exercise 1.12. Let
for n 1; the series
fn (x) =
x 4n1
x 4n+1
+
(4n 1)! (4n + 1)!
fn converges to F (x) = sin x x.
Solution of exercise 1.13

i) The power series for f (a + r e i ) may be integrated term by term because of its absolute
convergence in the disc centered at a of radius r < R. Since we have
Z 2
1 if k = n,
e i (kn) d = kn =
0 otherwise,
0
the stated formula follows.
49

2
ii) Similarly, expand f (a + r e i ) as a product of two series and integrate term by term.
Most contributions cancel out using the formula above, and only the terms |cn |2 r 2n
remain.
iii) If f is bounded on C, we have |cn r n | k f k . Letting r +, it follows that cn = 0
for n 1, which means that f is constant.

iv) The function sin is not bounded on C! Indeed, we have for instance lim sin(i x) =
x+
+. So there is no trouble.
Solution of problem 1
i)
(a) Let x I . For any p, q N, we have

f p (x) fq (x) sup f p ( y) fq ( y) = f p fq ,
yI
and this proves that the sequence

converges.
fn (x)
nN
is a Cauchy sequence in R, so it

(b) Let > 0 be fixed. There exists N such that f p fq for all p > q > N .
Let x I . We then have

f p (x) fq (x)
for any p > q > N ,
and since this holds for all p, we may fix q and let p . We obtain

f (x) fq (x)
for all q > N .
This bound holds independently of x I . Thus we have shown that

f fq
for any q N .
Finally, this being true for any > 0, it follows that the sequence ( fn )nN
converges uniformly to f .
Remark: At this point,
we havent proved that there is convergence in the normed

vector space E , kk . It remains to show that the limit f is in E , that is, that f is
continuous. This follows from Theorem 1.33, but we recall the proof.
(c) Let x I ; we now show that f is continuous at x.

Let > 0 be fixed. From the preceding question, there exists an integer N
such that k fn f k for all n N , and in particular k f N f k .
Since f N is an element of E , it is continuous. So there exists > 0 such that

y I
| y x| = f N ( y) f N (x) .
Using the triangle inequality, we deduce from this that for all y I such that
|x y| , we have

f ( y) f (x) f ( y) f N ( y) + f N ( y) f N (x) + f N (x) f (x) 3.
This proves the continuity of f at x, and since x is arbitrary, this proves that f
is continuous on I , and hence is an element of E .
(d) For any Cauchy sequence ( fn )nN in E , theprevious questions show that ( fn )nN
converges in E . Hence the space E , kk is complete.
ii) Let f be a fixed point of . Then we have

t I

f (t) = ( f ) (t) =
Moreover, it is easy to see that f (0) = ( f )(0) = 1.
t f (t)
1 + f (t)
2 .
50
Conversely, let f be a solution of the Cauchy problem (E ) + (C I ). Then ( f ) is

differentiable and we have
t f (t)
t I
( f ) (t) =
2 = f (t).
1 + f (t)
The functions f and ( f ) have the same derivative on I , and moreover satisfy
( f )(0) = 1
and
f (0) = 1.
It follows that ( f ) = f .

1 x2
we see that g (x) 1 for all x R. The
(1 + x 2 )2
mean-value theorem then proves that g is 1-Lipschitz, as stated.
iii) Looking at the derivative g : x 7

iv) Let f , g E . Then we have
Z
!

t

u
f
(u)
u
g(u)

( f ) (g) = sup
2
2 du

tI 0
1 + f (u)
1 + g(u)

Z t

f (u)
g(u)

sup
u
2
2 du

tI 0
1 + f (u)
1 + g(u)

Z a

f (u)
g(u)

u
2
2 du

1 + f (u)
1 + g(u)
0
by positivity. Using the inequality of the previous question, we get

Za
Za

a2
( f ) (g)
f (u) g(u) du k f gk
u
u du =
k f gk .
2
0
0
p
This is true for any f , g E , and hence is (a 2 /2)-Lipschitz; if 0 a < 2, this
map is a contraction.
v) According to the fixed-point theorem, the previous results show that has a unique
fixed point in E .
According to Question ii), this means that there existspa unique solution of the
Cauchy Problem (E ) + (C I ) on an interval [0, a] for a < 2.
To approximate the solution numerically, it is possible to select an arbitrary function
f0 (for instance, simply f0 = 0), and construct the sequence ( fn )nN defined by fn+1 =
( fn ) for n 0. This requires computing (numerically) some integrals, which is a fairly
straightforward matter (numerical integration is usually numerically stable: errors do
not accumulate in general25). The speed of convergence of the sequence ( fn )nN to the
solution f of the Cauchy problem is exponential: with I = [0, 1], the distance (from
the norm on E ) between fn and f is divided by 2 (at least) after each iterative step. It
is therefore possible to expect a good numerical approximation after few iterations (the
precision after ten steps is of the order of k f0 f k /1000 since 210 = 1024).
25
On the other hand, numerical differentiation tends to be much more delicate.
Chapter
Measure theory and the

Lebesgue integral
This chapter is devoted to the theory of integration. After a survey of the Riemann
integral and a discussion of its limitations, we present the Lebesgue integral. This requires a brief discussion of basic measure theory, which is important for probability
theory as well.
Note that the main theorems and techniques of integral calculus (change of variable, justification of exchanges of limits and integrals) are discussed in the next
chapter.
2.1
The integral according to Mr. Riemann
2.1.a
Riemann sums
One possible way of defining the integral of a function f on a finite interval

[a, b] is the following: start by subdividing the interval in n parts, which are
more or less of the same size (b a)/n, by choosing real numbers
a = x0 < x1 < x2 < < xn = b ;
52
a = x0 x1 x2
xk xk+1
xn1 xn = b
Fig. 2.1 The interval [a, b ] is here partitioned uniformly, with n subintervals of constant
length (b a)/n.
then approximate the function f by a step function which is constant on

each interval of the subdivision, and takes there the same value as f at the
beginning of the interval:
The sum of the areas of the small rectangles is then equal to
Sn =
n
X
k=1
(xk xk1 ) f (xk1 ),
and it may be expected (or hoped) that as n goes to infinity, this value will
converge to some limit. Riemanns result is that this happens if the approximations used for f improve steadily as n increases. This means that the values
taken by f on an interval [x, x + ] must be very close to f (x) when is
small. This is a continuity assumption for f . Here is the precise result.
DEFINITION 2.1 (Subdivision) A subdivision of the interval [a, b] is any tu-
ple of real numbers of the type

n = (x i ) i=0,...,n ,
where
a = x0 < x1 < x2 < < xn = b.
The step of the subdivision n is the length of the largest subinterval:

() = max x i+1 x i .
i[[0,n1]]
A marked subdivision of [a, b] is a pair

n = (x i ) i=0,...,n , (k )k=1,...,n ,

where (x i ) i=0,...,n is a subdivision of [a, b] and k xk1 , xk for all k.

DEFINITION 2.2 (Riemann sums) Let n = (x i ) i , (k )k be a marked subdivision of [a, b] and let f : [a, b] K be a function with values in K = R
The integral according to Mr. Riemann
53
or C. The Riemann sum of f associated to n is

S ( f , n ) =
n
X
(xk xk1 ) f (k ).
k=1
The following result provides a rigorous definition of the integral of a

continuous function on a finite interval:
THEOREM 2.3 (and definition of the integral) Let f : [a, b] K be a contin-
uous function. For any sequence (n )nN of marked subdivisions of [a, b], such that
the step of the subdivisions tends to 0 as n tends to infinity, the sequence of associated
Riemann sums converges. Moreover, the limit of this sequence is independent of the chosen sequence of marked subdivisions. The (Riemann) integral1 of the function f is
defined to be this common limit:
Zb
Zb
f (x) dx =
f = lim S ( f , n )
a
where (n )n satisfies lim (n ) = 0.

n
This result and definition may easily be extended to piecewise continuous

functions,2 and to even larger classes of functions, by means of a small change
in the construction used. Namely, define a step function on [ a, b] to be a
function which is constant on the subintervals of a subdivision
= (a0 , . . . , an )
of [a, b]. The
integral
of a step function associated to and taking the
value y i on a i1 , a i is obviously equal to
Zb
n
X
=
(a i a i1 ) y i .
a
i=1
For any function f : [a, b] R, define

Z

; step function with f
I ( f ) = sup
R
The notation d for differentiation, dy/dx, (which originally represented an s,
the initial of summa), are due to Gottfried Leibniz (16461716), who also invented the name
function and popularized the use of the equal sign = and of a dot to represent multipliRb
cation. Joseph Fourier (17681830) invented the notation a .
2
Recall that f : [a, b ] K is piecewise continuous if there exist an integer n and real
numbers
a = a0 < a1 < < an = b

such that, for all i [[0, n 1]], the function f restricted to a i , a i+1 is continuous and has
finite limits as x tends to a i from above, and to a i+1 from below.
1
54
and
I+ ( f ) = inf
Z
; step function with f
Then f is integrable in the sense of Riemann, or Riemann-integrable, if

we have I ( f ) = I+ ( f ).
A more delicate issue is to extend the definition to functions which are
defined on an arbitrary interval (not necessarily closed or bounded), such as
], ] or [, +[.
2.1.b Limitations of Riemanns definition

Some easily defined functions cannot be approximated by step functions, and the Riemann sums for those functions do not converge
to a common value. This is the case, for instance, for the function
d : [0, 1] R, which is defined by d (x) = 0 if x is irrational, and
d (x) = 1 if x is rational. This function is therefore not Riemannintegrable.
The definition turns out to be the very devil to deal with as soon as
questions of convergence of sequences of integrals appear. In particular,
to justify that the integral of the sum of a series is the sum of the series
of integrals, one needs very stringent conditions.
It also turns out that if we want to deal with a complete space of funcRb
tions with the norm a | f (t)| dt, not even the most liberal definition
of Riemann integral is sufficient. The smallest complete space containing C ([a, b]) is much larger than the space of Riemann-integrable
functions; it is the space L1 defined in the next section.
2.2
The integral according to Mr. Lebesgue
To solve the problems related to Riemanns definition of integrable function, the idea is to provide an alternate definition for which a much larger
class of functions will be integrable.3 In particular, this space of integrable
functions, denoted L1 , will turn out to be complete.
3
From the point of view of trying to obtain a complete space with the norm

Rb
k f k1 = a f (x) dx,
it is always possible to complete a metric space,

in particular, a normed vector space; more
precisely, given a normed vector space E , kk , there exists a complete normed vector space

F , kk such that E is isometric to a dense subspace E of F . Using this isometry, E may
55
f
xk+1
xk
Ak(n)
Fig. 2.2 Construction of the sets Ak(n) = f 1

value of Ak(n) k/2n .
2.2.a
xk , xk+1

, where xk = k/2n , with the
Principle of the method
Lebesgues idea to define the integral of a real-valued function is the following:

for various possible values (denoted ) of the function, we consider the size
of the set of points where f takes the value (or is suitably close to ). Multiply the size of this set by , and repeat with other values, tallying the total
sum thus obtained. Taking a limit in a suitable manner, this sum converges
to the integral of the function.4
In other words, instead of a vertical dissection as in Riemanns method,
one has to perform an horizontal dissection as follows. Assume first that
f takes non-negative values. Subdivide the set R+ of values ininfinitely
many intervals, say of size 21n . For each subinterval of the type 2kn , k+1
,
n
2
consider the set of real numbers x such that f (x) belongs to this interval (see
Figure 2.2):

k k+1
(n)
1
Ak = f
,
.
2n 2n

Denote by Ak(n) the size of this set (temporarily; this crucial notion
be seen as a dense subspace of a complete space. The general construction of F is somewhat

abstract (F is the set of Cauchy sequences in E , modulo the equivalence relation u v, if and
only if u v is a sequence converging to zero), but it is very powerful. In particular, it provides
one construction of the set R of real numbers, starting with Cauchy sequences in Q.
4
Henri Lebesgue used the following analogy: to sum up a certain amount of money, we can
instead of adding up the values of all coins as they come up one by one start by putting
together all dimes, all nickels, all $1 bills, and then count the number of each stack, multiply
it by the common value of each coin or bill in the stack, and finally add up everything.
56
Fig. 2.3 Construction of the integral of f .
of size will be made precise later), and then define

Sn ( f ) =
k
Ak(n) n .
2
k=0

It is quite easy to check that the sequence Sn ( f ) n is increasing (see Figure 2.4).
R
Hence it converges in R+ = R+ {+}, and the limit is denoted f .
The function f is said to be integrable (in the sense of Lebesgue) if the limit
is finite.5
This is the principle. As can be seen, the issue of evaluating the size of
the sets Ak(n) has been left open. This very important point is the subject of
the next sections.
Remark 2.4 Before going farther, we briefly explain what is gained from this new method:
the set of integrable functions is vast, beyond human imagination;6
the set of Lebesgue-integrable functions (not yet precisely defined),

denoted L1 or, to be
R
precise, L1 (R) or L1 [a, b ], is complete for the norm k f k = | f |;
the theorems of the Lebesgue integral calculus are infinitely simpler and more powerful
(in particular, concerning the inversion of limits and integrals, or differentiation under
the integral sign), even when applied to very regular functions.
5
The exact definition below will be slightly different, but equivalent; it will not involve a
particular choice of intervals (bounded by numbers of the form k/2n ) to subdivide R+ .
6
Much larger than the set of Riemann-integrable functions, in particular, because even very
irregular functions become integrable. It is true that some Riemann-integrable functions are
not Lebesgue integrable, and that this causes trouble, but they are typically functions which
create all the difficulties of the Riemann theory.
57
Fig. 2.4 Sn1 ( f ) (dark gray) is less than Sn ( f ) (light gray).
2.2.b Borel subsets

We now take up the question of measuring the size of a subset A of R. As
we will see, it turns out to be impossible to assign a measure to all subsets
of R. It is necessary to restrict which subsets are considered; they are called
measurable sets. Since it is natural to expect that the union A B of two
subsets which can be measured should also be measurable, and since analysis
quickly requires some infinite unions, the following definition is reasonable:
DEFINITION 2.5 Let X be an arbitrary set. The set of subsets of X is denoted
P(X ). A -algebra7 on X is any subset T P(X ) of the set of subsets of X
such that the following conditions hold:
(T1) X T and T ;
(T2) for all A T , the complement X \ A is also in T ;

S
(T3) if (An )nN is a sequence of elements of T , then
n=0 An T .
Hence a -algebra T is stable by complement and countable union. Moreover, since the complement of the union of a family of sets is the intersection
of their complements, T is also stable by countable intersections.
Example 2.6 P(X ) isa -algebra (the largest
-algebra on X ); {, X } is also a -algebra (the
smallest). For a X , , X , {a}, X \ {a} is a -algebra.
Consider a subset X of P(X ) (hence a set, the elements of which are

themselves subsets of X ). The intersection of all -algebras containing X is
still a -algebra, and of course it is contained in any -algebra containing X .
7
An algebra is a set of subsets stable under complement and finite union. The prefix
indicates the countable operations which are permitted.
58

mile Borel (18711956), was ranked first at the entrance examination of both the cole Polytechnique and the cole Normale
Suprieure. He chose the latter school and devoted himself to
mathematics; to Marguerite, daughter of Paul Appell (also a famous writer under the pseudonym Camille Marbo); and to politics (in the footsteps of Painlev). In 1941, he was jailed in
Fresnes because of his opposition to the regime of Vichy, and became a member of the resistance as soon as he was freed. In 1956,
he received the first Gold Medal from the CNRS. He founded
the Institut Henri-Poincar. Among his most important works,
we can mention that he studied the measurability of sets, and in
particular defined the sets of reals of measure zero and the Borel
sets, on which Lebesgue measure is well defined. He also worked
in probability theory and mathematical physics (the kinetic theory of gas). His student Henri Lebesgue (see page 61) used his
results to develop his theory of integration.
This is therefore the smallest -algebra containing X . This leads to the

following definition:
DEFINITION 2.7 Let X P(X ). The -algebra generated by X is the
smallest -algebra containing X . It is denoted T (X

X ). For A X , the
-algebra generated by A is the -algebra T (A) generated by X = {A}.

Example 2.8 For a X , the -algebra generated by {a} is , X , {a}, X \ {a} .
DEFINITION 2.9 (Borel sets) The Borel -algebra is the -algebra on R generated by the open subsets of R. It is denoted B(R). Its elements are called
Borel sets.
It is easy to check that the Borel -algebra is also the -algebra generated
by intervals of R of the type ], a], where a ranges over all real numbers.8
In particular, the following sets are Borel subsets of R:
1. all open sets;
2. all closed sets (since their complement is open);
3. all countable unions and intersections of open and closed sets.
However, there are many others! The Borel -algebra is much more complex
and resists any simple-minded exhaustive description.
8
This is of great importance in probability theory, where it is used to reduce the study of a
random variable to the study of its distribution function.
2.2.c
59
Lebesgue measure
It is unfortunately beyond the scope of this book to describe the construction

of the Lebesgue measure in all details. Interested readers are invited to look,
for instance, in Rudins book [76] or in Burks very clear course [16].
DEFINITION 2.10 A measurable space is a pair (X , T ), where X is an arbitrary set and T is a -algebra on X . Elements of T are called measurable
sets.
DEFINITION 2.11 Let (X , T ) be a measurable space. A measure on (X, T ),
or on X, by abuse of notation, is a map defined on the -algebra T , with

+ def
values in R = R+ {+}, such that
(MES1) () = 0;
(MES2) for any countable family (An )nN of pairwise disjoint measurable
sets in T , we have
X
[
Ai =
(Ai ).
i=0
n=0
The second property is called the -additivity of the measure.

A measure space9 is a triple (X , T , ), where (X , T ) is a measurable space
and is a measure on (X , T ).
def
Remark 2.12 Since a measure takes values in R+ = R+ {+}, it is necessary to specify the
addition rule involving +, namely,
x + = +
for all x R+ .
Note that the usual property (x + z = y + z) = (x = y) is no longer valid in full generality

(e.g., 2 + = 3 + ); however, it holds if z 6= +.
The following theorem, due to Lebesgue, shows that there exists a measure
on the Borel -algebra which extends the length of an interval.
THEOREM 2.13 (The Lebesgue measure on Borel sets) The pair (R, B) is a
measurable set and there exists on (R, B) a unique measure such that

[a, b] = |b a|
for all a, b R.
The idea for the construction of the Lebesgue measure is described in

Sidebar 1 on the following page.
The Lebesgue measure is translation-invariant (i.e., translating a measurable set does not change its measure), and is diffuse (that is, for any finite or
9
Note that very similar definitions will appear in Chapter 19, concerning probability spaces.
60
Sidebar 1 (Lebesgue measure for Borel sets)
Consider the length map which associates its length b a to any finite
interval [a, b]. We try to extend this definition to an arbitrary subset of R.
Let A P(R) be an arbitrary subset of R.
Let E (A) S
be the set of all sequences A = (An )nN of intervals such that A is
contained in
i=0 A i . For any A = (An )nN E (A), let
def
(A ) =
i=0
(Ai ),
and define the Lebesgue exterior measure of A as follows:

(A) =
def
inf (A ).
A E (A)
Unfortunately, this exterior measure is not a measure as defined previously: the

property of -additivity ( MES2) is not satisfied (this would contradict the result of
Sidebar 2).
On the other hand, one can show that the restriction of to the Borel -algebra
: B R
A 7 (A) = (A)
is a measure (this is by no means obvious or easy to prove). It is the Lebesgue
measure for Borel sets.
countable set , we have () = 0). In particular, (N) = (Z) = (Q) = 0.
Note that the converse is false: for instance, the Cantor triadic set is not countable, yet it is of measure zero (see Exercise 2.5 on page 68).
2.2.d The Lebesgue -algebra

The Lebesgue measure on Borel sets is sufficient for many applications (in
particular, it is commonly used in probability theory).
Nevertheless, it is interesting to know that the construction may be refined.
Consider a Borel set A with measure zero. There may exist subsets B A
which are not of measure zero, simply because B is not a Borel set, and hence
is not measurable. But it seems intuitively obvious that a subset of a set of
measure zero should also be considered negligible and be of measure zero (in
particular, measurable). In fact, it is possible to define a -algebra containing
all Borel sets, and all subsets of Borel sets with Lebesgue measure zero, and
to extend the Lebesgue measure to this -algebra. This -algebra is called the
Lebesgue -algebra and is denoted L . We have therefore B L . (The
construction of L and of the required extension of the Lebesgue measure are
again beyond the scope of this book.) We thus have:
61
Henri Lebesgue (18751941), son of a typographical worker,

entered the cole Normale Suprieure, where he followed the
courses of mile Borel. In 1901, he published in his thesis the
theory of the Lebesgue measure, generalizing previous work of
Borel (and also of Camille Jordan and Giuseppe Peano). The
following year, he developed his theory of integration, and used
it to study series of functions and Fourier series. He proved that
a bounded function is integrable in the sense of Riemann if and
only if it is continuous except on a negligible set. In his later
years, Lebesgue turned against analysis and would only discuss
elementary geometry.
PROPOSITION 2.14 Let A be an element of the Lebesgue -algebra such that
(A) = 0. Then any subset B A is in the Lebesgue -algebra and satisfies

(B) = 0.
This property is called completeness of the measure space (R, L , ) or of

the Lebesgue -algebra.10 The Borel -algebra is not complete, which justifies
the introduction of the larger Lebesgue -algebra.
Remark 2.15 In the remainder of this and the next chapter, we will consider the Lebesgue measure
on the Lebesgue -algebra, and a measurable set will be an element of L . Note that in practice
(for a physicist) there is really no difference between the Lebesgue and Borel -algebras.
2.2.e
Negligible sets
DEFINITION 2.16 A subset N X of a measure space (X , T , ) is negligi-
ble if it is contained in a measurable set N T such that (N ) = 0.

In particular, a subset N of R is negligible if and only if it is in the
Lebesgue -algebra and is of Lebesgue measure zero.
Remark 2.17 In terms of the simpler Borel -algebra, a subset N of R is negligible if and only
if there exists a Borel set N containing N which is of Lebesgue measure zero. But N itself is
not necessarily a Borel set.
Example 2.18 The set Q of rational numbers is negligible since, using -additivity, we have
(Q) =
X
rQ
X
{r } =
0 = 0.
rQ
DEFINITION 2.19 A property P defined for elements of a measure space

(X , T , ) is true for almost all x X or holds almost everywhere, abbreviated holds a.e., if it is true for all x X except those in a negligible
set.
10
This is not the same notion as completeness of a topological space.
62
Example 2.20 The Heaviside function
H : x 7
is differentiable almost everywhere.
The Dirichlet function
D : x 7
is zero almost everywhere.
1
0
0 if x < 0,
1 if x 0,
if x Q,
if x R \ Q,
DEFINITION 2.21 Let (X , T ) and (Y , T ) be measurable spaces. A function
f : X Y is measurable if, for any measurable subset B Y , the set

def
f 1 (B) = x X ; f (x) B
of X is measurable, that is, if f 1 (B) T for all B T .
This definition should be compared with the definition of continuity of a

function on a topological space (see Definition A.3 on page 574).
Remark 2.22 In practice, all functions defined on R or Rd that a physicist (or almost any
mathematician) has to deal with are measurable, and it is usually pointless to bother checking
this. Why is that the case? The reason is that it is difficult to exhibit a nonmeasurable function
(or even a single nonmeasurable set); it is necessary to invoke at some point the Axiom of Choice
(see Sidebar 2), which is not physically likely to happen. The famous Banach-Tarski paradox
shows that nonmeasurable sets may have very weird properties.11
2.2.f
Lebesgue measure on R n
Part of the power of measure theory as defined by Lebesgue is that it extends

fairly easily to product spaces, and in particular to Rn . First, one can consider
on R R the -algebra generated by sets of the type A B, where (A , B)
B(R)2 . This -algebra is denoted B(R) B(R); in fact, it can be shown
to coincide with the Borel -algebra12 of R2 with the usual topology. It is
11
Consider a solid ball B of radius 1 in R3 . There exists a way of partitioning the ball in
five disjoint pieces B = B1 B5 such that, after translating and rotating the pieces Bi
suitably without dilation of any kind the first two and the last three make up two distinct
balls of the same radius, not missing a single point:
7
It seems that this result flies in the face of the principle of conservation of volume during a
translation or a rotation. But that is the point: the pieces Bi (or at least some of them) are not
measurable, and so the notion of volume of Bi makes no sense; see [74].
12
In the sense of the -algebra generated by open sets.
63
then possible to prove that there exists a unique measure denoted =

on B(R) B(R) which is given on rectangles by the natural definition
(A B) = (A) (B)
for A, B B(R). This measure is called the product measure, or the

Lebesgue measure on R2 . The construction may be repeated, yielding a Borel
-algebra and a Lebesgue measure on Rn for all n 1.
2.2.g
Definition of the Lebesgue integral
Here we give a general definition of the Lebesgue integral on an arbitrary

measure space (X , T , ), which in a first reading the reader may assume
is R with the Lebesgue measure. It yields the same results as the construction
sketched in Section 2.2.a, but the added generality is convenient to prove the
important results of the Lebesgue integral calculus.
The first step is to define the integral for a non-negative measurable simple
function (i.e., a function taking only finitely many values). This definition
is then extended to non-negative functions, and finally to (some) real- or
complex-valued measurable functions.
DEFINITION 2.23 (Simple function) Let X be an arbitrary set, A X a subset of X . The characteristic function of A is the function A : X {0, 1}
such that A (x) = 1 if x A and A (x) = 0 otherwise.
Let (X , T ) be a measurable space. A function f : X R is a simple
function if there exist finitely many disjoint measurable sets A1 , . . . , An in T
and elements 1 , . . . , n in R such that
f =
n
P
i=1
i A .
i
Now if (X , T , ) is a measure space and f is a non-negative simple function

on X , the integral of f with respect to the measure is defined by
Z
n
X
def
f d =
i (Ai ),
i=1
which is either a non-negative real number or +. In this definition, the

convention 0 = 0 is used if either i = or (Ai ) = . If the integral
of f is finite, then f is integrable.
Remark 2.24 In the case where X = R and is the Lebesgue measure, the following notation
is also used for the integral of f with respect to :
Z
Z
Z
Z +
f d =
f d = f (x) dx =
f (x) dx.
R
64
Example 2.25 Let D be the Dirichlet function, which is 1 for rational numbers and 0 for
irrational numbers. Then D is a simple function (since Q is measurable), and

Z
Z
(1 D) d = 1 whereas
(1 D) d = +.
R
[0,1]
The definition is extended now to any measurable non-negative function.

Let (X , T , ) be a
measure space and let f : X R be a measurable function with nonnegative values. The integral of f is defined by
Z
Z

def
f d = sup
g d ; g simple such that 0 g f .
DEFINITION 2.26 (Integral of a non-negative function)
+
It is either a non-negative real number or +. If it is finite, then f is

integrable.
Finally, to deal with real-valued measurable functions f , the decomposition
of f into its non-negative and non-positive parts is used: The non-negative
def
part of a function f on X is the function f + = max( f , 0), and the nondef
positive part of f is the function f = max( f , 0). We have f = f + f
+
and the functions f and f are non-negative measurable functions.

f
f+
DEFINITION 2.27 (Integrability and integral of a function) Let (X , T , ) be
a measure space and let f : X 7 R be a measurable function, E X a

measurable set.R Then f is integrable
on E with respect to or f is R
integrable, if E f + d and E f d are both finite, or equivalently if the
non-negative function | f | is integrable on E . The integral of f on E with
respect to is then
Z
Z
Z
def
f d =
f + d
f d.
E
If f is a measurable complex-valued function, one looks separately at the real

and imaginary parts Re( f ) and Im( f ). If both are integrable on X with
respect to , then so is f and we have
Z
Z
Z
f d = Re( f ) d + i Im( f ) d.
65
In practice, the definitions above are justified by the fact that any measurable function f is the pointwise limit of a sequence ( fn )nN of simple
functions. If, moreover, f is bounded, it is a uniform limit of simple functions.13
As it should, the Lebesgue integral enjoys the same formal properties as
the Riemann integral. In particular, the following properties hold:
PROPOSITION 2.28 Let (X , T , ) be a measure space, and let f , g be measurable
functions on X , with values in R or C. Then we have:
i) if f and g are integrable with respect to , then for any complex numbers ,
C, the function f + g is integrable and
R
R
R
( f + g) d = f d + g d;
ii) R
if f and g are integrable and real-valued, and if f g, then
g d;
R
R
iii) if f is integrable, then f d | f | d;
f d
iv) R
if g is non-negative
and integrable and if | f | g, then f is integrable and
R
| f | d g d.
Finally, under some suitable assumptions,14 for any two measure spaces
(X , T , ) and (Y , T , ), the product measure space (X Y , T T , )
is defined in such a way that Fubinis formula holds: for any integrable
function f on X Y , we have
ZZ
Z Z

f d( ) =
f (x, y) d( y) d(x)
X Y
Z X Z Y

=
f (x, y) d(x) d( y)
Y
(an implicit fact here is that the integral with respect to one variable is a
measurable function of the other variable, so that all integrals which appear
make sense). In other words, the integral may be computed in any order (first
over X then Y and conversely). Moreover, a measurable function f defined
on X Y is integrable if either

Z Z
Z Z

f (x, y) d( y) d(x) or

f (x, y) d(x) d( y)

X
13
This should be compared with the Riemann definition of integral, where step functions
are used. It turns out that there are relatively few functions which are uniform limits of step
functions (because step functions, constant on intervals, are much more rigid than simple
functions can be, defined as they are using arbitrary measurable sets).
14
Always satisfied in applications.
66
is finite.
All this applies in particular to integrals on R2 with respect to the Lebesgue
measure , and by induction on Rd for d 1. The Fubini formula
becomes
ZZ
Z Z

f (x, y) dx dy =
f (x, y) dy dx
R2
R
R
Z Z

=
f (x, y) dx dy.
R
The following notation is also commonly used for integrals on R2 :

ZZ
ZZ
ZZ
f d2 =
f d( ) =
f (x, y) dx dy
R2
R2
R2
ZZ
=
f (x, y) d(x) d( y).
R2
2.2.h Functions zero almost everywhere,

space L1 of integrable functions
PROPOSITION 2.29 Let (X , T , ) be a measure space, for instance, X = R with
the Lebesgue -algebra and measure. Let f : X R+ {+} be a non-negative

measurable function. Then we have
Z
f d = 0 f is zero almost everywhere.
Proof. See Exercise 2.7 on page 69.
In particular, if the values of an integrable function f are changed on a

set of measure zero, the integral of f does not change. We have
Z
Z
f = g a.e. =
f d = g d.
In the chapter concerning distributions, we will see that two functions
which differ only on a set of measure zero should not be considered to be
really different. Hence it is common to say that such functions are equal,
instead of merely15 equal almost everywhere.
DEFINITION 2.30 (L1 spaces) Let (X , T , ) be a measure space. The space
L1 (X, ) or simply L1 (X) if no confusion concerning can arise, is the
15
This is not simply an abuse of notation. Mathematically, what is done is take the set of
integrable functions, and take its quotient modulo the equivalence relation being equal almost
everywhere. This means that an element in L1 (X ) is in fact an equivalence class of functions.
This notion will come back during the course.
67
vector space of -integrable complex-valued functions on X , up to equality

almost everywhere.
In particular, L1 (R
R ) or simply L1 denotes the vector space of functions
integrable on R with respect to Lebesgue measure, up to equality almost everywhere.
For any interval [a, b] of R, L1 [ a, b] is the space of functions on [a, b]
integrable with respect to the restriction of Lebesgue measure to [a, b]. This
integrability condition can also be expressed as saying that the function on R
obtained by extending f by 0 outside [a, b] is in L1 (R).
The following important theorem justifies all the previous theory:
THEOREM 2.31 (Riesz-Fischer) Let (X , T , ) be a measure space. ThenR the space
L1 (X , ) of complex-valued integrable functions, with the norm k f k1 = | f | d,

is a complete normed vector space. In particular, the spaces L1 (R) and L1 [a, b] are
complete normed vector spaces.
Moreover, the space C ([a, b]) of continuous functions on [a, b] is dense in the
space L1 [a, b], which is therefore the completion of C ([a, b]), and similarly for the
space C (R) of continuous functions on R and L1 (R).
2.2.i And today?

Lebesgues theory of integration became very quickly the standard approach
at the beginning of the twentieth century. However, some other constructions
remain useful, such as that of Stieltjes16 (see Sidebar 8). Generalizations of the
integral, such as the Ito stochastic integral,17 or recently the work of Ccile
DeWitt-Morette and Pierre Cartier on the Feynman path integrals [53], have
also appeared. There is also a pedagogical18 definition of an integral which
combines advantages of the Lebesgue and Riemann integrals, due to Kurzweil
and Henstock (see, e.g., [93]).
16
The so-called Riemann-Stieltjes integral defines integration on R with respect to a Borel

measure; they can be used to avoid appealing to distributions in some situations, and are still
of great use in probabiliy theory.
17
This concerns Brownian motion in particular.
18
This is not the only reason it is interesting.
68
EXERCISES
Exercise 2.1 Let d N . Prove or disprove (with a counterexample) the following state-
ments:
i) An open set in Rd is bounded if and only if its Lebesgue measure is finite.

ii) A Borel set in Rd has positive Lebesgue measure if and only if it contains a non-empty
open set.
iii) For any integer p and 0 < p < d , the set R p {0}d p has Lebesgue measure zero in Rd .
Exercise 2.2 Find the -algebras on R generated by a single element.
Exercise 2.3 Show that each of the following subsets generate the Borel -algebra B(R):
1. closed sets in R;
2. intervals of the type ]a, b ];
3. intervals of the type ], b ].
Show that B(Rd ) is generated by
1. closed sets in Rd ;
2. bricks, that is, sets of the type {x Rd | a i < xi b i , i = 1, . . . , d };
3. closed half-spaces {x Rd | xi b }, where 1 i d and b R.
Exercise 2.4 Let (X , T , ) be a measure space. Prove the following properties:
i) () = 0;
ii) is increasing: if A , B T and A B, then we have (A) (B);
iii) let (An )nN be an increasing sequence of measurable sets in T , and let A =
Then we have (An ) (A);
n
An .
iv) let (An )nN beT

a decreasing sequence of measurable sets in T such that (A1 ) < +,
and let A = n An . Then we have (An ) (A). Show also that this is not
n
necessarily true if (A1 ) = +;
v) let (An )nN be any sequence of measurable sets in T . Then we have

[ X
An
(An )
(-subadditivity of the measure).
n
Exercise 2.5 In this exercise, we denote by O the set of open subsets in R for the usual
topology, and (as usual) by B and L the Borel -algebra and the Lebesgue -algebra. We have
O B L P(R), and we will show that all inclusions are strict.
i) Show that O $ B.
ii) Let C be the Cantor triadic set, defined by

\
P
def
C=
Fp ,
where
F p = x R ; (kn )nN E p , x = kn /3n ,
p1
n1
n
o
E p = (kn )nN {0, 1, 2}N ; n p, kn 6= 1 .
def
a) Show that C is a compact set with (C ) = 0.
Exercises
69
b) Show that C is in one-to-one correspondance with {0, 1}N , and in one-to-one correspondence with R.
c) Show that L is in one-to-one correspondance with P(R).
d) One can show that there exists an injective map from B to R. Deduce that the
second inclusion is strict.
iii) Use Sidebar 2 on the following page to conclude that the last inclusion is also strict.
Exercise 2.6 Let f be a map from R to R. Is it true that the two statements below are
equivalent?
i) f is almost everywhere continuous;

ii) there exists a continuous function g such that f = g almost everywhere.
Exercise 2.7 Let f be a measurable function
on a measure space (X , T , ), such that f (x) >
R
0 for x A with (A) > 0. Show that f d > 0.

Deduce from this a proof of Theorem 2.29 on page 66.
Exercise 2.8 (Egorovs theorem) Let (X , T , ) be a finite measure space, that is, a measure
space such that (X ) < + (for instance, X = [a, b ] with the Lebesgue measure). Let
( fn )nN be a sequence of complex-valued measurable functions on X converging pointwise to a
function f .
For n N and k N , let
\

1
def
.
x X ; f p (x) f (x)
En(k) =
k
pn

For fixed k, show that the sequence En(k) n1 is increasing (for inclusion) and that the measure
ofEn(k) tends to (X ) as n tends to infinity. Deduce the following theorem:
THEOREM 2.32 (Egorov) Let (X , T , ) be a finite measure space. Let ( fn )nN be a sequence of
complex-valued measurable functions on X converging pointwise to a function f . Then, for any > 0,
there exists a measurable set A with measure (A ) < such that ( fn )nN converges uniformly to f on
the complement X \ A .
This theorem shows that, on a finite measure space, pointwise convergence is not so
different from uniform convergence: one can find a set of measure arbitrarily close to the
measure of X such that the convergence is uniform on this set.
70
Sidebar 2 (A nonmeasurable set) We describe here the construction of a subset
of R which is not measurable (for the Lebesgue -algebra).

Consider the following relation on R: for x, y R, we say that x y if and
only if (x y) Q.
This relation is of course an equivalence relation.a
We now appeal to the Axiom of Choice of set theoryb to define a subset E [0, 1]
containing exactly one element in each equivalence class of c. Note that the use
of the Axiom of Choice means that this set E is not defined constructively: we know
it exists, but it is impossible to give an explicit definition of it.
We now claim that E is not in the Lebesgue -algebra.
To prove this, for any r Q, let Er = {x + r ; x E } .
If r, s Q and r 6= s, the sets Er and E s are disjoint. Indeed, if there exists
x Er E s , then by definition there exist y, z E such that x = y + r = z + s,
which implies that z y = r s is a rational number, and so z y. But E
contains a single element in each equivalence class, so this means that y = z and
hence r = s.
Assuming now that E is measurable, let = (E ). Then for each r Q, Er
is measurable with measure (Er ) = (because it is a translate of E ).
[
def
Let
S=
Er .
r Q[1,1]
As a countable union of measurable sets, this set is itself measurable. If 6= 0, then

(S ) = + since S is a disjoint union, so we can use -additivity to compute
its measure. But clearly S [1, 2] so (S ) 3. The assumption > 0 is
therefore untenable, and we must have = 0. Then (S ) = 0. But we can also
easily check that
[0, 1] S [1, 2] .
This implies that (S ) 1, contradicting this last assumption.

It follows that E is not measurable.
a
This means that, for all x, y, z :

x x;
x y if and only if y x;
if x y and y z, then x z.
The Axiom of Choice states for any non-empty set A, there exists a map f : P(A)
A, called a choice function, such that f (X ) X for any non-empty subset X A. In
other words, the choice function selects, in an arbitrary way, an element in any non-empty
subset of X . It is interesting to know that a more restricted form of the Axiom of Choice
(the Axiom of Countable Choice) is compatible with the assertion that any subset of R is
Lebesgue-measurable. (But is not sufficient to prove this statement). The interested reader
may look in books such as [51] for a very intuitive approach, or [45] which is very clear.
c
This means that for any x R, there is a unique y E with x y.
71
SOLUTIONS
i) Wrong. Take, for instance, d = 1 and the unbounded open set
[
1
U =
k, k + k
2
k1
P
which has Lebesgue measure (U ) = k1 2k = 1.
ii) Wrong. For d = 1, consider for example the Borel set (R\Q)[0, 1], which has measure
1 but does not contain any non-empty open set.
iii) True. It is easy to describe for any > 0 a subset of Rd containing R p {0}d p with
measure less than .
Solution of exercise 2.6. Wrong. Let f = D be the Dirichlet function, equal to 1 on Q

and zero on R \ Q. Then f is nowhere continuous, but f = 0 almost everywhere.
Solution of exercise 2.7. For any n 1, let

def
En = x X ; f (x) > 1/n .
The sequence (En )nN is an increasing sequence of measurable sets, and its union is

x X ; f (x) > 0 ,
which has positive measure since it contains A. Using the -additivity, there exists n0 such
that En0 has positive measure. Then the integral of f is at least equal to (En0 )/n0 > 0.
Solution of exercise 2.8. For any k N, we have En(k)

(X ). It is then possible
n

to construct a strictly increasing sequence (nk )kN such that X En(k)
< /2k for all k N,
k
and then conclude.
Chapter
Integral calculus
3.1
Integrability in practice
A physicist is much more likely to encounter concrete functions than
abstract functions for which integrability might depend on subtle theoretical
arguments. Hence it is important to remember the usual tricks that can be
used to prove that a function is integrable.
The standard method involves two steps:
first, find some general criteria proving that some standard functions
are integrable;
second, prove and use comparison theorems, which show how to reduce
the integrability of a complicated function to the integrability of one
of the standard ones.
3.1.a
Standard functions
The situation we consider is that of functions f defined on a half-open interval [a, b[, where b may be either a real number or +, and where the
functions involved are assumed to be integrable on any interval [a, c ] where
a < c is a real number (for instance, f may be a continuous function on
[a, b[). This is a fairly general situation, and integrability on R may be reduced to integrability on [0, +[ and ], 0].
Integral calculus
74
PROPOSITION 3.1 (Integrability of a non-negative function) Let a be a real
and let b R {}. Let f : [a, b[ R+ be a non-negative measurable

function which has a primitive F . Then f is integrable on [a, b[ if and only if F has
a finite limit in b.
Example 3.2 The function log is integrable on ]0, 1]. The functions x 7 e ax are integrable
on [0, +[ if and only if a > 0.
PROPOSITION 3.3 (Riemann functions t ) The function f : t 7 t with
R is integrable on [1, +[ if and only if < 1. It is integrable on ]0, 1] if and

only if > 1.
PROPOSITION 3.4 (Bertrand functions t log t) The function t 7 t (log t)
with and real numbers is integrable on [1, +[ if and only if

< 1
or
3.1.b
( = 1
and < 1).
Comparison theorems
Once we know some integrable functions, the following simple result gives
information concerning the integrability of many more functions.
PROPOSITION 3.5 (Comparison theorem) Let g : [a, b[ C be an integrable
function, and let f : [a, b[ C be a measurable function. Then
i) if | f | |g|, then f is integrable on [a, b[;

ii) if f = O(g) or if f = o(g) and if f is integrable on any interval [a, c ] with
b
c < b, then f is integrable on [a, b[;

iii) if g is non-negative, f g, and f is integrable on any interval [a, c ] with
b
c < b, then f is integrable on [a, b[.
Exercise 3.1 Study the integrability of the following functions (in increasing order of diffi-
culty):
i) t 7 cos(t) log(tan t) on ]0, /2[;
ii) t 7 (sin t)/t on ]0, +[;
iii) t 7 1/(1 + t sin2 t) on [0, +[, where > 0.
PROPOSITION 3.6 (Asymptotic comparison) Let g be a non-negative measurable function on [a, b[, and let f be a measurable function on [a, b[. Assume that
g is not integrable and that f and g are both integrable on any interval [a, c ] with
c < b. Then asymptotic comparison relations between f and g extend to their integrals
as follows:
Rx
Rx
Rx
Rx
f g = a f a g
and
f = o(g) = a f = o
g .
a
b
xb
xb
Exchanging integrals and limits or series
75
If g and h are integrable on [a, b[, then asymptotic comparison relations between
f and h extend to the remainders as follows:
R b
Rb
Rb
Rb
f h = x f x h
and
f = o(h) = x f = o
h .
x
b
xb
Exercise 3.2 Show that
+
1
xb
e x t
dt log x.
p
1 + t 2 x0+
3.2
Exchanging integrals and limits or series
Even when dealing with concrete functions, Lebesgues theory is extremely useful because of the power of its general statements about exchanging
limits and integrals. The main theorem is Lebesgues dominated convergence
theorem.
THEOREM 3.7 (Lebesgues dominated convergence theorem) Let (X , T , )
be a measure space, for instance, R or Rd with the Borel -algebra and the Lebesgue
measure. Let ( fn )nN be a sequence of complex-valued measurable functions on X . Assume that ( fn )nN converges pointwise almost everywhere to a function f : X R,
and moreover assume that there exists a -integrable function g : X R+ such that
| fn | g almost everywhere on X , for all n N . Then
i) f is -integrable;
ii) for any measurable subset A X , we have
Z
Z
fn d =
f d.
lim
n
This result is commonly used both to prove the integrability of a function

obtained by a limit process, and as a way to compute the integral of such
functions.
1
n
Example 3.8 Suppose we want to compute lim
able y = nx, we are led to

In =
ey
2 /n2
n e x cos x
dx. By the change of vari1 + n2 x 2
cos( y/n)
dy.
1 + y2
Integral calculus
76
The sequence of functions fn : y 7 e y /n cos( y/n)/(1 + y 2 ) converges pointwise on R to

the function f : x 7 1/(1 + x 2 ). Moreover, | fn | f , which is integrable on R (by comparison
with t 7 1/t 2 at ).
Therefore, by Lebesgues theorem, the limit exists and is equal to
Z
+
1
1 + 1
dx =
arctan x = 1.
1 + x 2
Another useful result does not require the domination assumption, but
instead is restricted to increasing sequences. This condition is more restrictive,
but often very easy to check, in particular for series of non-negative terms.
THEOREM 3.9 (Beppo Levis monotone convergence theorem) Let (X , T , )
be a measure space, for instance, R or Rd with the Borel -algebra and the Lebesgue
measure. Let ( fn )n be a sequence
of non-negative measurable functions on X . Assume that the sequences fn (x) n are increasing for all x X . Then the function
f : X R {+}, defined by
for all x X ,
f (x) = lim fn (x)

n
is measurable and we have

lim
fn d =
(this quantity may be either finite or infinite).
f d
This result is often used as a first step to prove that a function is integrable.
Here finally is the version of the dominated convergence theorem adapted
to a series of functions:
THEOREM 3.10 (Term by term integration of a series) Let (X , T , ) be a mea-
sure P
space, for instance, R or Rd with the Borel -algebra and the Lebesgue measure.
Let
fn be a series of complex-valued measurable functions on X such that
Z
X
| fn | d < +.
n=0
Then the series of functions

fn converges absolutely almost everywhere on X and
its sum is -integrable. Moreover, its integral may be computed by term by term
integration: we have
Z X
Z
X
fn d =
fn d.
n=0
n=0
Integrals with parameters
77
3.3
Integrals with parameters
3.3.a
Continuity of functions defined by integrals
Let (X , T ; ) be a measure space. Assume we have a function of two variables

f : [a, b] X R or C. Assuming the integrals make sense, one can define
a function by integrating over X for fixed values of the parameter x [a, b]:
let
Z
def
I (x) =
f (x, t) dt
for x [a, b] .
X
A natural question is to determine what properties of regularity the integral I (x) will inherit from properties of f .
The two main results concern the continuity or derivability of f . Here
again, the Lebesgue theory gives very convenient answers:
THEOREM 3.11 (Continuity of functions defined by integrals) Let x0 [a, b].
Assume that for almost all t X , the function x 7 f (x, t) is continuous at x0 , and
that there exists a -integrable function g : X R such that

f (x, t) g(t) for almost all t X ,
R
for all x in a neighborhood of x0 . Then x 7 X f (x, t) dt is continuous at x0 .
Proof. This is a simple consequence of the dominated convergence theorem. Indeed,

let (xn )nN be any sequence of real numbers converging to x0 . For n large enough, the
values of xn lie in the neighborhood of x0 given by the statement of the theorem. Then
for almost all t X , we have lim f (xn , t) = f (x0 , t) by continuity, and, moreover,
n

f (xn , t) g(t) for all t. From the dominated convergence theorem, we deduce
Z
Z
Z
lim I (xn ) = lim
f (xn , t) dt =
f (x0 , t) dt =
lim f (xn , t) dt, = I (x0 ),
n
which, being valid for an arbitrary sequence (xn )nN , is the statement we wanted to
prove.
Counterexample 3.12 It is of course natural to expect that the conclusion of Theorem 3.11
should be true in most cases. Here is an example where, in the absence of the domination
property, the continuity of a function defined by an integral with parameters does not hold.
Let f be the function on R R defined by
1 2
exp
t
if x 6= 0,
2
x
f (x, t) =
0
if x = 0.
For all t R, the

R function x 7 f (x, t) is continuous on R, and in particular at 0.
Let I (x) = f (x, t) dt.
Integral calculus
78
For any fixed nonzero x, the function t 7 f (x, t) is a gaussian centered at 1/x, so its
integral is easy to compute; it is in fact equal to 1. But for x = 0, the function t 7 f (0, t) is
identically zero, so its integral is also zero. Hence we have
I (0) = 0
x R ,
but
I (x) = 1.
This shows that the function I is not continuous at zero.
3.3.b
Differentiating under the integral sign
THEOREM 3.13 (Differentiability of functions defined by integrals) Let x0
[a, b]. Assume that there exists a neighborhood V of x0 such that for x V , the
function x 7 f (x, t) is differentiable at x0 for almost all t X , and t 7 f (x, t)
is integrable. Assume, moreover, that the derivatives are dominated by an integrable
function g, that is, there exists an integrable function g : X R such that

f

(x, t) g(t)
for all x V and almost all t X .
x

R
Then the function x 7 I (x) = f (x, t) dt is differentiable at x0 , and moreover we
have
Z

Z

f
d

f (x, t) dt
=
(x0 , t) dt.
dx
X
X x
x=x0
Those theorems are very useful in particular when dealing with the theory
of the Fourier transform.
Counterexample 3.14 Here again, the assumption of domination cannot be dispensed with.
For instance, let

F (x) =
cos(x t)
dt;
1 + t2
one can then show (see Exercise 3.10 on page 84) that
F (x) = e |x|
2
for all x R. This function is continuous on R, but is not differentiable at 0.
3.3.c
Case of parameters appearing in the integration range
THEOREM 3.15 With X = R and the same assumption as in Theorem 3.13, let
moreover v : V R be a function differentiable at x0 . Let
Z v(x)
I (x) =
f (x, t) dt.
0
Then I is differentiable at x0 and we have

Z v(x)

Z v(x)

d
f

f (x, t) dt
=
(x0 , t) dt + v (x) f x0 , v(x0 ) .

dx
x
0
0
x=x0
Double and multiple integrals
79
The derivative of this function defined by an integral involves two terms:

one obtained by keeping the limits of integration constant, and differentiating
under the integral sign; and the second obtained by differentiating the integral
with varying boundary.
v(x)
0
f
(x , t) dt
x 0
f (x0 + h, )
Cnt h 2
f (x0 , v(x0 ))
h v (x) f x0 , v(x0 )
f (x0 , )
v(x0 )
v(x0 + h)
h v (x0 )
3.4
Double and multiple integrals
For completeness and ease of reference, we recall the theorem of Fubini
mentioned at the end of the previous chapter.
THEOREM 3.16 (Fubini-Lebesgue) Let (X , T , ) and (X , T , ) be measure
spaces, each of which is the union of a sequence of sets with finite measure, for instance,
X = Y = R with the Lebesgue measure. Let f : X Y R be a mesurable
function on X Y for the product -algebra. Then:
i) f is integrable on X Y with respect to the product measure if and only
if one of the following integrals is finite:
Z Z

Z Z

f (x, y) dy dx
f (x, y) dx dy.
or
ii) If f is integrable on X Y , then for almost all x X , the function y 7

f (x, y) is -integrable on Y , and similarly when exchanging the role of x
and y.
iii) If f is integrable on X Y , the function defined for almost all x X by
Z
def
x 7 I (x) =
f (x, y) d( y)
Y
Integral calculus
80
A child prodigy, the Venetian Guido Fubini (18791943) found,

in nineth grade, a formula for as a series converging faster
than anything known at the time. At the Scuola normale superiore di Pisa, he was student of Dini (see page 22). He then held
positions at Pisa and Torino. Mussolinis fascist regime forbade
him to teach, and he went into exile to Paris and then to Princeton University in the United States, where he finished his career.
In addition to his work on the Lebesgue integral, he proved important results in functional analysis, differential geometry, and
projective geometry.
is -integrable on X and we have

Z
ZZ
I (x) d(x) =
X
X Y
f d( )(x, y).
Remark 3.17 Since X and Y play symmetric roles, we also have
ZZ
f =
X Y
Z Z
Y

Z Z

f (x, y) d(x) d( y) =
f (x, y) d( y) d(x).
X
For non-negative functions, we have a converse:

THEOREM 3.18 (Fubini-Tonelli) Let (X , T , ) and (X , T , ) be measure spaces,
each of which is the union of a sequence of sets with finite measure, for instance,
X = Y = R with the Lebesgue measure. Let f : X Y R be a non-negative
measurable function. Then the integral of f on X Y with respect to the product
measure , which is either a non-negative real number or +, is given by the
Fubini formula
Z Z

ZZ
f (x, y) d( )(x, y) =
f (x, y) d(x) d( y)
X Y
Y
X
Z Z

=
f (x, y) d( y) d(x).
X
In particular, if the following two conditions hold:

i) for almost all x X , the function y 7 f (x, y) is integrable on Y ;
def R
ii) the function defined for almost all x X by x 7 I (x) = R f (x, y) dy is
-integrable,
then f is integrable on X Y with respect to the product measure .
Change of variables
81
3.5
Change of variables
The last important point of integral calculus is the change of variable formula in multiple integrals. Suppose we wish to perform a change of variable
in 3 dimensions
(x, y, z) 7 (u, v, w) = (x, y, z),
(and analogously in n dimensions); then the jacobian of this change of variable
is defined as follows:
DEFINITION 3.19 Let and be two open sets in R3 . Let be a bijection
from to , which is continuous and has continuous differential, as well

as its inverse. Denote (u, v, w) = (x, y, z). The jacobian matrix of is the
matrix of partial derivatives
u u u
x y z
v
def D(u, v, w) def v
.
=
Jac (x, y, z) =
D(x, y, z) x y z
w w w
x y z
This is a continuous function on . The jacobian of the change of variable

is the determinant of the jacobian matrix:
J (x, y, z) = det Jac (x, y, z).
It is also a continuous function on .

The generalization of this definition to change of variables in Rn is immediate.
THEOREM 3.20 (Change of variables) Let U and V be open sets in Rn . Let
: U V be a bijection which is continuous, has continuous differential, and is

such that the jacobian of does not vanish on U . Then for any function f : V R p ,
we have
f is integrable on V
and moreover
Z
Z
f (u1 , . . . , un ) du1 . . . dun =
V
( f ) J is integrable on U

f (x1 , . . . , xn ) J (x1 , . . . , xn ) dx1 . . . dxn .
Here integration and integrability are with respect to the Lebesgue measure on Rn , and
mean that each component of f is integrable.
Integral calculus
82
As an example, take the change of variable in polar coordinates in the

plane R2 , denoted : (r, ) 7 (x, y), which we may assume to be defined on
the open set ]0, +[ ] , [. We can write

r
=
x 2 + y2,
x = r cos ,
or
y
y = r sin ,
p
.
= 2 arctan
x + x 2 + y2
(Tricky question: why shouldnt we write = arctan( y/x) as is often seen?
Why do most people write it anyway?)
The jacobian of this change of variable is

x y

D(x, y) r r cos
sin
=
=

D(r, ) x y r sin r cos = r,

and we recover the well-known formula:
ZZ
Z Z
f (x, y) dx dy =
f (r, ) r dr d.
Remark 3.21 We might be interested in the inverse transformation. The computation of
D(r,)
D(x, y)
directly from the formula is rather involved, but a classical result of differential calculus avoids
this computation: we have
D(r, ) D(x, y)
= 1,
D(x, y) D(r, )
and this result may be generalized to arbitrary locally bijective differentiable maps in any
dimension.
Remark 3.22 In one dimension, the change of variable formula is
f (u) du =

f (x) | (x)| dx.
It is important not to forget the absolute value in the change of variable, for instance,
Z
Z
Z
f (x) dx =
f ( y) |1| dy =
f ( y) dy.
R
Sometimes absolute values are omitted, and one writes, for instance, dy = dx. In this case,
the range of integration is considered to be oriented, and the orientation is reversed when the
change of variable is performed if it corresponds to a decreasing function : x 7 y = (x).
Thus we have
Z +
Z
Z +
f (x) dx =
f ( y) (dy) =
f ( y) dy,
with the minus sign cancelling the effect of interchanging the bounds of integration.
Exercises
83
EXERCISES
|sin t|
2
dt
log x.
x+
t
Exercise 3.4 Check that the jacobian for the change of variables from cartesian coordinates
to spherical coordinates (x, y, z) 7 (r, , ) in R3 is J = r 2 sin .
Exercise 3.5 Find (and justify!) the following limits:
lim
1
0
n /2 x
dx
1 + n2 x 2
and
lim
1/
|sin x| n
dx.
1 + x2
Exercise 3.6 Let f () be the distribution function of energy emitted by a black body, as a
function of the frequency. Recall that, by definition, the total amount of energy with frequency
between 0 and 1 is given by the integral
Z 1
f () d.
0
Max Planck showed at the end of the nineteenth century that the distribution function f ()
is given by the formula (which now bears his name)
f () =
A 3
,
} 1
e 2h
where A is a constant, = 1/kb T is the inverse temperature, and h} is the Planck constant.
def
What is the distribution function g() of the energy as function of the wavelength = c /?
For which frequency max is f () maximal? What is the wavelength for which g() is maximal?
Why is it that max 6= c /max ?
Exercise 3.7 Show that the function x 7 arctan(x) arctan(x 1) is integrable on R (with
respect to Lebesgue measure, of course) and show that
Z +

arctan(x) arctan(x 1) dx = .
Generalize this formula.
Exercise 3.8 (Counterexample to Fubinis theorem) For (x, y) [0, 1]2 , define
f (x, y) =
R1
sgn( y x)
.
max(x, y)2
f (x, y) dx for all y [0, 1].
R 1 R 1
Deduce from this the value of 0 0 f (x, y) dx dy.
R 1 R 1
ii) Similarly, compute 0 0 f (x, y) dy dx.
i) Compute
iii) Any comments?
Integral calculus
84
Exercise 3.9 (Integrability of a radial function) Let f be a (measurable) function on R3 , such
that we have
f (x, y, z) r
r0
p
def
2
2
for some real number , where r = x + y + z 2 . Assuming that f is integrable on any
bounded region in R3 {0}, discuss the integrability of f at the origin, depending on the
value of . Generalize this to functions on Rd .
Exercise 3.10 (Fourier transform of a Lorentzian function) Let
F (x) =
cos u x
du
1 + u2
for all x R.
i) Where is f defined? Is it continuous on its domain of definition?

Z +
x cos t
ii) Let g(x) =
dt. What is the relation between F and g?
2
2
x + t
iii) Deduce from the previous question a differential equation satisfied by g. (Start by
comparing
2
x
2
x
and
.
x2 x2 + t2
t2 x2 + t2
Solve the differential equation, and deduce the value of F .)
85
SOLUTIONS
Solution of exercise 3.4. Let (r, , ) denote spherical coordinates (where is latitude and
longitude). We have
x = r cos cos ,
and therefore the jacobian is

cos cos

J = r sin sin
r cos sin
y = r cos sin ,
cos sin
r sin sin
r cos cos
z = r sin ,
sin
r cos
0

= r 2 sin .

Solution of exercise 3.6. The function g() has the property that the amount of energy
radiated by frequencies in the interval [0 , 1 ], corresponding to wavelengths 0 and 1 , is
Z 0
equal to
g() d. The change of variable = c / leads to
1

Z 0
Z 1
Z 1
d
c
g() d =
g(c /) d =
g(c /) 2 d,
d
1
0
0
and this shows that
g(c /)
.
2
Note that there is no reason for the functions 7 f () and 7 2 f () to reach a maximum
at the same point. In fact, rather than the functions f and g themselves, it is the differential
forms (or measures) f () d and g() d which have physical meaning.
f () = c
Solution of exercise 3.7. Integrability comes from the inequalities

1
1
arctan(x) arctan(x 1)
,
2
1 + (x 1)
1 + x2
which are consequences of the mean value theorem. To compute the value of the integral, let
F denote a primitive of arctan. Then we have
ZM

arctan(x) arctan(x 1) dx = F (M ) F (M 1) F (M ) + F (M 1)
M
= arctan( M ) arctan( M
),
where M [M 1, M ] and M
[M 1, M ] are given by the mean value theorem
again. The result stated follows immediately.
Solution of exercise 3.10.

1
i) The uniform upper bound u 7 1+u
2 and the continuity theorem show that F is defined
and continuous on R. It is obvious that F (x) , with equality for x = 0.
ii) By the change of variable t = x y, we have for x > 0
Z +
cos(x y)
g(x) =
dy = F (x).
2
1 + y
It follows that g(x) has a limit as x tends to zero on the right, namely,
g(0+ ) = F (0) = .
Similarly, for x < 0 we have g(x) = F (x) and hence g(0 ) = . Since the definition
gives g(0) = 0, we see that the function g is not continuous at 0.
86
Integral calculus
iii) A simple computation shows that the two derivatives mentioned in the statement of the
exercise are equal. From this, we deduce that
Z +
Z +
2
x
2
x
2
g(x)
=
cos
t
dt
=
cos t 2
dt
2
2
2
2
2
2
x
x
x +t
t
x +t
Z +
x
=
cos t
dt,
x2 + t2
using two integrations by parts. It follows that g is a solution of the differential equation
g (x) g(x) = 0.
The general solution of this equation is of the form
g(x) = e x + e x ,
where (+ , + ) are real constants (which give the value of g on R+ ), and ( , )
similarly on R .
Since g is bounded, the coefficient + of e t on R+ is zero, and the coefficient of
t
e on R also. Since g is odd and g(0+ ) = , we finally find that
g(x) = sgn(x) e |x| .
Since F (x) = sgn(x) g(x) and F (0) = , we finally find that
x R
F (x) = e |x| .
Chapter
Complex Analysis I
This chapter deals with the theory of complex-valued functions of one complex
variable. It introduces the notion of a holomorphic function, which is a function
f : C defined and differentiable on an open subset of C. We will see that
the (weak) assumption of differentiability implies, in strong contrast with the real
case, the (much stronger) consequence that f is infinitely differentiable. We will
then study functions with singularities at isolated points and which are holomorphic
except at these points; we will see that their study has important applications to
the computation of many integrals and sums, notably for the computation of the
Fourier transforms that occur in physics. In Chapter 5, we will see how techniques
of conformal analysis provide elegant solutions for certain problems of physics in
two dimensions, in particular, problems of electrostatics and of incompressible fluid
mechanics, but also in the theory of diffusion and in particle physics (see also
Chapter 15).
4.1
Holomorphic functions
Whereas differentiability in R is a relatively weak constraint,1 we will see
during the course of this chapter that differentiability in terms of a complex
variable implies by contrast many properties and rigidifies the situation, in
a sense that will soon be made precise.
1
A function f , defined on an open interval I of R and with values in R or C may be
differentiable at all points of I without, for instance, the derivative being continuous. For
instance, if we set f (0) = 0 and f (x) = x 2 sin(1/x) for all nonzero x, then f is differentiable
at all points of R, f (0) = 0, but f is not continuous at 0 (there is no limit of f at 0).
Complex Analysis I
88
The life of Baron Augustin-Louis Cauchy (17891857) is quite

eventful since Cauchy, for political reasons (as a supporter of the
French monarchy), went into exile in Italy and then Prague. He
was a long-time teacher at the cole Polytechnique. His work is
extremely wide-ranging, but it is to analysis that Cauchy brought
a fundamental contribution: a new emphasis on rigorous arguments. He was one of the first to worry about proving the convergence of various processes (in integral calculus and in the theory
of differential equations); he worked on determinants and introduced the notation of matrices with double indices. In complex
analysis, his contributions are fundamental, as shown by the many
theorems bearing his name.
4.1.a
Definitions
Let us start by defining differentiability on the complex plane. Let be an

open subset of the complex2 plane C, and let f : C be a complex-valued
function of a complex variable. Let z0 be given.
DEFINITION 4.1 The function f : C is said to be differentiable at
z0 C if
f (z) f (z0 )
z z0
exists, i.e., if there exists a complex number, denoted f (z0 ) such that

f (z) f (z0 ) f (z0 ) (z z0 ) = o(z z0 ) [z z0 ].
f (z0 ) = lim
def
zz0
One may ask whether this definition is not simply equivalent to differentiability in R2 , after identifying the set of complex numbers and the real plane
(see sidebar 3 on page 134 for a reminder on basic facts about differentiable
functions on R2 ). The following theorem, due to Cauchy, shows that this is
not the case and that C-differentiability is stronger.
THEOREM 4.2 (Cauchy-Riemann equations) Let f : C be a complexvalued function of one complex variable; denote by fe the complex-valued function
on R2 naturally associated to f , that is, the function
fe :
R2 C,
def
(x, y) 7 fe(x, y) = f (x + i y).
While we are speaking of this, we can recall that the study of the complex plane owes
much to the works of Rafaele Bombelli (15261573) (yes, so early!), who used i to solve algebraic
equations (he called it di meno, that is, [root] of minus [one]), and of Jean-Robert Argand
(17681822) who gave the interpretation of C as a geometric plane. The complex plane is
sometimes called the Argand plane. It is also Argand who introduced
the modulus for a
p
complex number. The notation i, replacing the older notation 1, is due to Euler (see
page 39).
89
Then the function f is C-differentiable at the point z0 = x0 + i y0 if and only if
fe is R2 -differentiable
Cauchy-Riemann equation.
fe
fe
(x0 , y0 ) = i (x0 , y0 ).
x
y
Proof
Assume that f is differentiable in the complex sense at a point z0 . We can define
a linear form d fe(z0 ) : R2 C in the following manner:
def

d fe(z0 ) (k, l ) = f (z0 ) k + il .
We have then, according to the definition of differentiability of f at z0 ,

e
f (x, y) fe(x0 , y0 ) d fe(x0 , y0 ) (x x0 ), ( y y0 ) = o(z z0 ),
which shows that fe is R2 -differentiable at the point (x0 , y0 ). Moreover, since
it follows that
d fe(x0 , y0 )(k, l ) = f (z0 )k + i f (z0 )l ,
d f (z0 ) = f (z0 ) dx + i f (z0 ) dy.

Comparing with the formula
(see page 134), we get
d fe =
f
= f (z0 )
x
fe
fe
dx +
dy
x
y
and
f
= i f (z0 ),
y
and therefore the required equality.

Conversely: if fe is R2 -differentiable and satisfies fe/ x = i fe/ y, then

fe

fe
fe
dx +
d y (k, l ) =
dx + i dy (k, l )
d fe(x0 , y0 ) (k, l ) =
x
y
x
=
fe
(k + il ),
x
which shows, by the definition of differentiability of fe, that

f (z) f (z0 ) f (z z0 ) = o(z z0 ),

x
and thus that f is differentiable and that its derivative is equal to f (z0 ) =
fe
(z ).
x 0
We henceforth write f instead of fe, except in special circumstances.

In the following, we will very often write a complex-valued function in the
form
f = P +iQ
with P = Re( f ), Q = Im( f ).
The relation obtained in the preceding theorem allows us to write
Complex Analysis I
90
THEOREM 4.3 (Cauchy-Riemann equations) Let f : C be a function dif-
ferentiable at z0 . Then we have the

P
Q
=
x (x0 , y0 )
y (x0 , y0 )
following equations:
Cauchy-Riemann equations.

Q
P
=
x (x0 , y0 )
y (x0 , y0 )
DEFINITION 4.4 Let be an open subset of C. A function f : C is
holomorphic3 at z0 C if it is differentiable at z0 , and it is holomorphic

on if it is holomorphic at every point of . We will denote by H () the
set of all functions holomorphic on .
If A is a non-empty subset of C, a function f defined on A is holomorphic on A if there exist an open neighborhood of A and a holomorphic
function g on such that g(z) = f (z) for z A.
The following exercise (which is elementary) is an example of the rigidity
of holomorphic functions that we mentioned before.
DEFINITION 4.5 A domain of C is a connected 4 open subset of C.
Exercise 4.1 Show that if f : C is a function holomorphic on a domain of C, the
following properties are equivalent:

(i) f is constant on ;
(ii) P = Re( f ) is constant on ;

(iii) Q = Im( f ) is constant on ;
(iv) | f | is constant on ;
(v) f is holomorphic on .
(Solution on page 129)
4.1.b Examples
Example 4.6 The functions z 7
z 2,
z 7 e z , the polynomials in z, are holomorphic on C.
Example 4.7 The function z 7 1/z is holomorphic on C .
On the other hand, the function f : z 7 z is not differentiable at any

point. Indeed, let us show, for instance, that f is not differentiable at 0. For
3
The name holomorphic was introduced by Jean-Claude Bouquet (18191885) and his colleague Charles Briot (18171882). Some texts use equivalently the words analytic or regular.
4
A topological space is said to be connected if it is not the union of two disjoint open sets;
in other words, it is a set which is in one piece; see Appendix A.
91
this, we compute
f (z) f (0)
z
= .
z 0
z
If now we let z follow a path included in R approaching 0, this quotient is
equal to 1 thoughout the path, whereas if we let z follow a path included in
the set of purely imaginary numbers, we find a constant value equal to 1;
this shows that the quantity of interest has no limit as z 0, and thus that f is
not differentiable at 0. If we add a constant, we see that it is not differentiable
at z0 for any z0 .
Example 4.8 The function z 7 |z| is nowhere holomorphic, and similarly for the functions
z 7 Re(z) and z 7 Im(z).
As a general rule, holomorphic functions will be those which depend on

z only (and of course sufficiently smooth), whereas those functions that
dependpnot only on z but also on z will not be holomorphic, as for instance
|z| = z z or Re(z) = z+z
. One may say that the variables z and z are
2
treated as if they were independent variables.5
4.1.c
The operators / z and / z
We consider the plane R2 , which we will freely identify with C.

DEFINITION 4.9 We will denote by dx the linear form which associates its
first coordinate to any vector (k, l) in R2 . Similarly, we denote by dy the

linear form which associates its second coordinate to any vector (k, l) in R2 :
dx(k, l) = k
def
dy(k, l) = l.
def
We will write, moreover, dz = dx + i dy and dz = dx i dy.

Exercise 4.2 Show that the differential form dz = dx + idy is the identity (dz = IdC ) and
that the differential form dz is the complex conjugation.
For any function f : C, R2 -differentiable (and therefore not necessarily differentiable in the complex sense), we have (check sidebar 3 on page 134):
df =
f
f
dx +
dy.
x
y
The functions z = x + i y and z = x i y are R2 -differentiable on R2 C;

moreover, dz = dx +i dy and dz = dx i dy. We deduce that dx = 12 (dz +dz )
5
For those readers interested in theoretical physics, this treatment as independent variables
can be found also in field theory, for instance in the case of Grassmann variables.
Complex Analysis I
92
and dy =
1
(dz
2i
dz ), which allows us to write:

if
1f
(dz + dz ) +
(dz dz)
2 x
2 y

1 f
f
1 f
f
=
i
dz +
+i
dz .
2 x
y
2 x
y
df =
This leads us to the following definition:

DEFINITION 4.10 We introduce the differential operators
def 1
=
z
2
i
x
y
def 1
=
z
2
and
+i
x
y
We will sometimes denote them and , when there is no ambiguity.

With these operators, the differential of a function R2 -differentiable can therefore be written in the form
df =
f
f
dz +
dz .
z
z
We can then write another form of Theorem 4.2:

THEOREM 4.11 (Cauchy-Riemann equations) If f is C-differentiable, that is,
holomorphic, at z0 , we have
f
(z ) = f (z0 )
z 0
and
f
(z ) = 0.
z 0
This last relation is none other than the Cauchy-Riemann condition. One may write
also the two preceding equations as
f
f
f
(z ) =
(z ) = i
(z ).
z 0
x 0
y 0
Example 4.12 The function f : z 7 Re(z) =
point.
z+z
2
is not holomorphic because
f
z
1
2
at every
Remark 4.13 We could also change our point of view and look at functions that can be
expressed in terms of z only:
DEFINITION 4.14 A function f : C is antiholomorphic if it satisfies
f
=0
z
for all point in .
Cauchys theorem
93
Remark 4.15 Be careful not to argue that f / z = 0 at all points, so the function
is constant, which is absurd. The function z 7 z does satisfy the condition, and
therefore is antiholomorphic, but it is certainly not constant! The quantity f / z is
equal to f (z) only for a holomorphic function. For an arbitrary function, it does
not have such a clear meaning.
4.2
Cauchys theorem
In this section, we will describe various versions of a fundamental result of
complex analysis, which says that, under very general conditions, the integral
of a holomorphic function around a closed path is equal to zero. We will start by
making precise what a path integral is, and then present a simple version of
Cauchys theorem (with proof), followed by a more general version (proof
omitted).
4.2.a
Path integration
We now need a few definitions: argument, path, and winding number of a path
around a point.
DEFINITION 4.16 (Argument) Let be an open subset of C = C \ {0} and
let : R be a continuous function such that for z we have

z
e i(z) = .
|z|
Such a function is called a continuous determination of the argument. If

is connected, then two continuous determinations of the argument differ
by a constant integral multiple of 2. There exists a unique continuous
determination of the argument
: C \ R ], [
which satisfies (x) = 0 for any real number x > 0. It is called the principal
determination of the argument and is denoted Arg.
Remark 4.17 One should be aware that a continuous determination of the argument
does not necessarily exist for a given open subset. In particular, one cannot find one
for the open subset C itself, since starting from argument 0 on the positive real axis
and making a complete circle around the origin, one comes back to this half-axis with
argument equal to 2, so that continuity of the argument is not possible.
DEFINITION 4.18 (Curve, path) A curve in an open subset in the complex
plane is a continuous map from [0, 1] into (one may also allow other
Complex Analysis I
94
sets of definitions and take an arbitrary real segment [a, b] with a < b). If
: [0, 1] is a curve, we will denote by e its image in the plane; for
simplicity, we will also often write simply without risk of confusion. A
curve is closed if (0) = (1). It is simple if the image doesnt intersect
itself, which means that (s) = (t) only if s = t or, in the case of a closed
curve, if s = 0, t = 1 (or the opposite). A path is a curve which is piecewise
continuously differentiable.
DEFINITION 4.19 (Equivalent paths) Two paths
: [a, b] C
: [a , b ] C
and
are called equivalent if their images coincide and the corresponding curves
have the same orientation: ([a, b]) = ([a , b ]), with (a) = (a ),
(b) = (b). In other words, there exists a re-parameterization of the path
, preserving the orientation, which transforms it into , which means a
bijection u : [a, b] [a , b ] which is continuous, piecewise differentiable
and strictly increasing, and such that (x) = u(x) for all x [a, b].
Example 4.20 Let :
[0, 2] C
and
[0, /2] C
t 7 e
t 7 e 2i cos t .
These define equivalent paths, with image equal to the unit circle.
it
If : [a, b] is a path, we have of course the formula

Zb
(b) (a) =
(t) dt.
a
On the other hand, if f : R is a real-valued function, twice continuously

differentiable, then denoting by 1 and 2 the real and imaginary parts of ,
we have

Z b
d1 f
d2

f
(t)
+
(t)
dt,
(4.1)
f (b) f (a) =
x
dt
y
dt
a
the right-hand side of which is nothing else, in two dimensions, than

Z
f d.
One must be aware that this formula holds only for a real-valued function.
DEFINITION 4.21 (Integral on a path) Let f : C be an arbitrary (mea-
surable) function (not necessarily holomorphic), and a path with image

contained in . The integral of f on the path is given, if the right-hand
side exists, by
Z
Zb

def
f (z) dz =
f (t) (t) dt.
(4.2)
The path is called the integration contour.
Cauchys theorem
95
One checks (using the change of variable formula) that this definition
does not depend on the parameterization of the path , that is, if is a path
equivalent to , the right-hand integral is simultaneously defined and gives
the same result with or . We then have the following theorem:
THEOREM 4.22 Let f be a holomorphic function on and : [a, b] C a path
with image contained in . Then

f (b) f (a) =
f
(z) dz.
z
(4.3)
Proof. We write f = f1 + i f2 and then for each of the real-valued functions f1

and f2 we appeal to formula (4.1). Formula (4.2), applied to f / z, and the Cauchy
conditions then imply the result.
Remark 4.23 This formula is only valid if f is a holomorphic function. If we take the
example of the function f : z 7 z , which is not holomlorphic, we have f / z = 0
at all points of C. Formula (4.3) does not apply, since otherwise it would follow that
f is identically zero! An extension of the formula, valid for any complex function
(holomorphic or not) is given in Theorem 4.51 on page 106 (the Green formula).
DEFINITION 4.24 The length of a path is the quantity

def
L =
b
a

(t) dt.
We then have the following intuitive property: the integral of a function

along a path is at most the length of the path multiplied by the supremum
of the modulus of the function.
PROPOSITION 4.25 If f is holomorphic on and if is a path contained in ,
then
Z

f (z) dz sup f (z) L .

z
4.2.b Integrals along a circle

Some integrals can be computed easily using a suitable parameterization of
the path; this is the case for instance for integrals along circles.
Consider the circle C (a ; R) centered at the point a with radius R, with positive (counterclockwise) orientation; one can use the parameterization given by
: [0, 2] C,
7 a + R e i ,
and, therefore, we have z = a + R e i and dz = iR e i .
Complex Analysis I
96
PROPOSITION 4.26 (Integral along a circle) Let a C and R > 0. Let f be a

function defined on an open subset in C containing the circle = C (a ; R). Then
Z
Z 2
f (z) dz =
f (a + Re i ) iR e i d.
0
As an example, let us compute the integral of the function z 7 1/z on

the unit circle. The path can be represented by : [0, 2] C, 7 z = e i .
We find
Z
Z 2 i
dz
ie
=
d = 2i.
(4.4)
e i
z
0
4.2.c Winding number

THEOREM 4.27 (and definition) Let be a closed path in C, U the complement
of the image of in the complex plane. The winding number of around the
point z is given by
Z
1
d
def
Ind (z) =
for all z U .
(4.5)
2i z
Then Ind (z) is an integer, which is constant as a function of z on each connected

component of U and is equal to zero on the unbounded connected component of U .
Proof. See Appendix D.
Example 4.28 A small picture is more eloquent than a long speech.
The picture on the right shows the various values of the function Ind for
a curve in the complex plane. One sees that the meaning of the value of the
function Ind is very simple: for any point not on , it gives the number
of times the curve turns around z, this number being counted algebraically
(taking orientation into account). Being in the unbounded connected component of U means simply being entirely outside the loop.
1
2
0
1
2
0
Exercise 4.3 Show by a direct computation that if is a circle with center a and radius r ,
positively oriented, then

Ind (z) =
1
0
if |z a| < r ,
if |z a| > r .
4.2.d Various forms of Cauchys theorem

We are now going to show that if f H (), where is open and has no
holes, and if is a path in , then the integral of f along is zero.
We start with the following simple result:
PROPOSITION 4.29 Let be an
R open set and a closed path in . If F
H () and F is continuous, then
F (z) dz = 0.
Cauchys theorem
Proof. Indeed,
(b ) = (a).
97

F (z) dz = F (b ) F (a) = 0 since is a closed path so
In particular we deduce the following result.

Z
any closed path and any n N

For or
, we have z n dz = 0.
any closed path , 0
/ and n = 2, 3, . . .
THEOREM
4.30
Proof. For any n Z, n 6= 1, the function z 7 z n is the derivative of

z 7
z n+1
,
(n + 1)
which is holomorphic on C if n 0 or on C \ {0} if n 2. The preceding

proposition gives the conclusion.
This property is very important, and will be used to prove a number of

results later on. An important point: the (excluded) case n = 1 is very
special. We have shown (see formula (4.4)) that the integral of 1/z along at
least one closed path is not zero. This comes from the fact that the function
z 7 1/z is the only one among the functions z 7 1/z n which is not the
derivative of a function holomorphic on C . It is true that we will define a
logarithm function in the complex plane, with derivative equal to z 7 1/z,
but we will see (Section 5.1.a on page 135) that this complex logarithm can
never be defined on the whole of C , but only on smaller domains where it is
not possible for a closed path to turn completely around 0.
Here is now a first version, somewhat restricted, of Cauchys theorem.
PROPOSITION 4.31 (Cauchys theorem for triangles) Let be a closed triangle in the plane (by this is meant a filled triangle, with boundary consisting
of three line segments),
entirely contained in an open set . Let p and let
f H \ { p} be a function continuous on and holomorphic on except
possibly at the point p. Then
Z
f (z) dz = 0.
Proof. The proof, elementary but clever, is well worthy of attention. Because of its
length, we have put it in Appendix D.
Remark 4.32 We have added some flexibility to the statement by allowing f not to be holomorphic at one point (at most) in the triangle. In fact, this extra generality is illusory. Indeed,
we will see further on that if a function is holomorphic in an open subset minus a single
point, but is continuous on the whole open set, then in fact it is holomorphic everywhere on
the open set.
The following version is much more useful.
Complex Analysis I
98
THEOREM 4.33 (Cauchys theorem for convex sets) Let be an open convex
subset of the plane and let f H () be a function holomorphic on . Then, for

any closed path contained in ,
Z
f (z) dz = 0.
Proof. Fix a point a and, for all z , define, using the convexity property,
Z
def
F (z) =
f ( ) d ,
[a,z]
where we have denoted

def
[a, z] = (1 )a + z ; [0, 1] .
We now show that F is a primitive of f .

def
Let z0 be a point in . For any z , the triangle = a, z, z0 is contained in .
R
Applying Cauchys Theorem for a triangle, we get that [z ,z] f (z) dz = F (z) F (z0 ),
0
hence
Z

F (z) F (z0 )
1
f (z0 ) =
f ( ) f (z0 ) dz 0
zz0
z z0
z z0 [z0 ,z]
using the continuity of f . The function F is therefore differentiable at z0 and its
derivative is given by F (z0 ) = f (z0 ). .
We have proved that f = F . Proposition 4.29 then ensures that the integral of f
on any closed path is zero.
Finally, there exists an even more general version of the theorem, which is
optimal in some sense. It requires the notion of a simply connected set (see Definition A.10 on page 575 in Appendix A). The proof of this is unfortunately
more delicate. However, the curious reader will be able to find it in any good
book devoted to complex analysis [43, 59].
THEOREM 4.34 (Cauchys theorem for simply connected sets) Let be a simply connected open subset of the plane and let f H (). Then for any closed
inside ,
Z
f (z) dz = 0.
Remark 4.35 In fact, one can show that if f H (), where is an open simply connected
set (as for instance a convex set), then f admits a primitive on . The vanishing of the integral
then follows as before.
In the case of an open set with holes, one can use a generalization of
Cauchys theorem. It is easy to see that such an open set can be dissected in a
manner similar to what is indicated in Figure 4.1. One takes the convention
that the boundary of the holes is given the negative (clockwise) orientation,
and the exterior boundary of the open set is given the positive orientation.
The integral of a holomorphic function along the boundary of the open set,
with this orientation, is zero.
Properties of holomorphic functions
99
Fig. 4.1 One can dissect an open set having one or more holes to make it simply
connected. The integral on the boundary of the new open set thus defined
is zero since is simply connected, which shows that the integral on (with
the orientation chosen as described on the left) is also zero.
4.2.e
Application
One consequence of Cauchys theorem is the following: one can move the
integration contour of an holomorphic function continuously without changing the value of the integral, under the condition at least that the function
be holomorphic at all the points which are swept by the path during its
deformation.
THEOREM 4.36 Let be an open subset of C, f HR(), and aRpath in . If
can be deformed continuously to a path in , then
f (z) dz =
f (z) dz.
Proof. The proof is somewhat delicate, but one can give a very intuitive graphical
description of the idea. Indeed, one moves from the contour in Figure 4.2 (a) to the
contour of Figure 4.2 (b) by adding the contour of Figure 4.2 (c). But the latter does
not contain any hole (since is deformed continuously to ); therefore the integral
of f on the last contour is zero.
4.3
4.3.a
The Cauchy formula and applications
THEOREM 4.37 (Cauchys formula in a convex set) Let be a closed path in a
convex open set and let f H (). If z and z

/ , we have
Z
1
f ( )
f (z) Ind (z) =
d .
2i z
Complex Analysis I
100
(a)
(b)
(c)
Fig. 4.2 The three contours , et .
In particular, if turns once around z counterclockwise, then

Z
1
f ( )
f (z) =
d .
2i z
Proof. Let z , z
/ and define a function g by
f ( ) f (z)
if , 6= z,
def
z
g( ) =
f (z)
if = z.
Then g H ( \ {z}) and is continuous at z; it is then possible to show (see Theorem
4.56) that g is holomorphic on . So, by Cauchys theorem for a convex set,
R
g( ) d = 0, hence
Z
Z
f ( )
d
0=
d +
f (z)
Z
f ( )
=
d + 2i Ind (z) f (z)
z
by definition of the winding number of a path around a point.
Remark 4.38 This result remains valid for a simply connected set.
Let us apply this theorem to the case, already mentioned, where the path
of integration is a circle centered at z0 with radius r: = C (0 ; r). One can
then parameterize by z = z0 + re i and we have
Z
Z
1
f (z)
1
f (z0 + re i ) i
f (z0 ) =
dz =
rie d
2i z z0
2i
r e i
Z
1
=
f (z0 + re i ) d,
2
which allows us to write:
THEOREM 4.39 (Mean value property) If is an open subset of C and if f is
holomorphic on , then for any z0 and any r R such that B(z0 ; r) , we

have
f (z0 ) =
1
2
101
f (z0 + re i ) d.
This theorem links the value at a point of a holomorphic function with

the average of its values on an arbitrary circle centered at this point, which
explains the name.
Here now is one the most fundamental results from the theoretical viewpoint; it also has practical implications, as we will see in Chapter 15 about
Green functions.
THEOREM 4.40 (Cauchy formulas) Let be an open subset of C and let f :
C be an arbitratry function. We have
f H () f is analytic on .
Moreover, for any n N and for any a ,

Z
n!
f (z)
f (n) (a) =
dz,
2i (z a)n+1
where is any path inside such that Ind (a) = 1.
Recall that f is analytic on an open subset if and only if, at any point
of , f admits locally a power series expansion converging in an open disc
with non-zero radius. The proof will use Cauchys Formula:
P
Proof. Let us prove the implication . If the power series f (z) =
n=0 c n (z
a)n converges on an open ball
r > 0 around the point a, then according
P with radius
to Theorem 1.69, the series ncn (z a)n1 converges on the same open ball. Without
loss of generality, we can work with a = 0. Denote
X
def
ncn z n1 ,
z B(0 ; r ).
g(z) =
n=0
We now show that we have indeed f = g on B(0 ; r ).

Let w B (0 ; r ). Choose a real numer such that |w| < < r . Let z 6= w be a
complex number. We have
n
X
f (z) f (w)
z wn
g(z) =
cn
nw n1 .
zw
zw
n=1
The term between brackets is zero for n = 1. For n 2, it is equal to

n
n1
X
z wn
nw n1 = (z w)
kw k1 z nk1 .
zw
k=1
(To prove this formula, it P

suffices to express the right-hand side using the well-known
k
z
w / z .) Hence, if |z| < ,
relation 1 x n = (1 x) n1
k=0 x , putting x = / w or x =
the modulus of the term between brackets is at most equal to n(n 1)n2 /2, and

X
f (z) f (w)

g(z) |z w|
n2 |cn | n2 .
zw
n=1
Complex Analysis I
102
But this last series is convergent (Theorem 1.69). Letting z tend to w, we find that f is
differentiable at w and that g(w) = f (w).
So it follows that f is holomorphic
Pon B(a ; r ), and its derivative admits a power
series expansion of the type f (z) = n1 ncn (z a)n1 . By induction one can even
show that the successive derivatives of f are given by the formula in Theorem 1.69 on
page 35, which is therefore valid for complex variables as well as for real variables.
Let us prove the implication . We assume then that f H ().
Let a , and let us show that f admits a power series expansion around a.
Since is open, one can find a real number R > 0 such that B(a ; R) . Let be
a circle centered at a and with radius r < R. We have (Theorem 4.37)
Z
1
f ( )
f (z) =
d
z B(a ; r ).
2i z
But, for any , we have
hence the geometric series

z a |z a|

< 1;
a
r
a
z
X
(z a)n
1
=
n+1
(
a)
z
n=0
converges uniformly with respect to on , and this for any z B(a ; r ), which allows us
to exchange the series and the integral and to write
Z X
Z
(z a)n
1 X
f ( )
1
n
f (z) =
f
(
)
d
=
(z
a)
d
2i n=0 ( a)n+1
2i n=0
( a)n+1
X
=
cn (z a)n ,
n=0
with
Z
1
f ( )
d .
2i
( a)n+1
So we have shown that f admits a power series expansion around a. Moreover we now
(see the general theory of power series) that in fact we have cn = f (n) (a)/n!, and this
gives the formula of the theorem.
In conclusion, f is indeed analytic on and the formula is proved.
def
cn =
An immediate consequence of this theorem is:

THEOREM 4.41 Any function holomorphic on an open set is infinitely differen-
tiable on .
We can now pause to review our trail so far. We have first defined a
holomorphic function on as a function differentiable at any point of .
Having a single derivative at any point turns out to be sufficient to prove
that, in fact, f has derivatives of all order, which is already a very strong
result. Even stronger is that f is analytic on , that is, that it can be expanded
in power series around every point. And the proof shows something even
beyond this: is we take a point z , and any open ball centered at z entirely
103
Fig. 4.3 If f is holomorphic on , then f has a power series expansion converging on

the gray open ball (the open ball with largest radius contained in ).
contained in , then f can be expanded in a power series that converges on

this ball6 (see Figure 4.3). In particular, we have the following theorem.
THEOREM 4.42 Let f be a function holomorphic on the whole of C. Then f admits
around any point a power series expansion that converges on the whole complex plane.
Example 4.43 The function z 7 exp(z) has, around 0, the power series expansion given by
ez =
X
zn
n!
n=0
with infinite radius of convergence. Around a point z0 C, the function has (another)
expansion
ez =
X
e z0 (z z0 )n
n=0
which converges also on the whole of C.
n!
DEFINITION 4.44 An entire function is any function which is holomorphic

on the whole complex plane.
The following theorem is due to Cauchy, but bears the name of Liouville7:
THEOREM 4.45 (Liouville) Any bounded entire function is constant.
This theorem tells us that any entire function which is not constant must
necessarily tend to infinity somewhere. This somewhere is of course at
infinity in the complex plane, but one must not think that it means that
f (z) tends to infinity when |z| tends to infinity. For instance, if we consider
the function z 7 exp(z), we see that if we let z tend to infinity on the left
6
There can of course be convergence even outside this ball; this brings the possibility of
analytic continuation (see Section 5.3 on page 144).
7
Not because it was thought that Cauchy already had far too many theorems to his name
but because, Liouville having quoted it during a lecture, a German mathematician attributed
the result to him by mistake; as sometimes happens, the mistake took hold.
Complex Analysis I
104
Joseph Liouville (18091882), son of a soldier and student at

the cole Polytechnique, became professor there at the age of 29,
a member of the Acadmie des Sciences at 30 and a professor
at the Collge de France at 42. A powerful and eclectic mind,
Liouville studied with equal ease geometry (surfaces with constant
curvature), number theory (constructing the first transcendental
number, which bears his name), mechanics (three-body problem,
celestial dynamics,...) analysis,... and he became actively involved
in the diffusion of mathematics, creating an important journal,
the Journal de mathmatiques pures et appliques, in 1836.
side of the plane (that is, by letting the real part of z go to ), then exp(z)
tends to 0. It is only if the real part of z tends to + that |e z | tends to
infinity.
Proof of Liouvilles theorem. We start by recalling the following result:
LEMMA 4.46 If f is an entire function such that f (z) 0 on C, then f is constant.
According to the Cauchy formulas (Theorem 4.40), if is a contour turning once

around z in the positive direction, we have
Z
1
f ( )
d = f (z).
2i ( z)2

If we assume that f is bounded, there exists M R such that f (z) < M for
any z C.
Let then z C, and let us show that f (z) = 0.
Let > 0. Consider the circle centered at z and with arbitrary radius R > 0.
We have
Z

Z

f ( )
M
f ( )
1

f (z) 1
d
d
.
2i ( z)2 2 | z|2
R
If we select the radius R so that R M /

(which
is always possible since the function
f is defined on C), we obtain therefore f (z) .
As a consequence, we have found that f (z) = 0. This being true for any z C, the
function f is constant.
4.3.b
Maximum modulus principle
THEOREM 4.47 (Maximum modulus principle) Let be a domain in C, f a
function holomorphic on , and a . Then, either f is constant, or f has no local

maximum at the point
a, which means that any neighborhood of a contains points b
such that f (a) < f (b).
What does this result mean? At any point in the open set of definition,
the modulus of the holomorphic function cannot have a local maximum. It
can have a minimum: it suffices to consider z 7 z which a local (in fact,
global) minimum at z = 0. But we can never find a maximum.
105
COROLLARY 4.48 Let f H () and let V be a bounded open subset of , such
that the closure V of V is contained in . Then the supremum of f on V is achieved

on the boundary of V , which means that there exists w V such that

f (w) = sup f (z).
zV
We will see later that this theorem also applies to real harmonic functions
(such as the electrostatic potential in a vacuum), which will have important
physical consequences (see Chapter 5; it is in fact there that we will give,
page 141, a graphical interpretation of the maximum modulus principle).
4.3.c Other theorems

We write down for information a theorem which is the converse of Cauchys
theorem for triangles:
THEOREM 4.49 (Morera) Let f be a complex-valued function, continuous on an
open subset of C, such that
f (z) dz = 0
for any closed triangle . Then f H ().

This theorem has mostly historical importance, and some interest for pure
mathematicians, but is of little use in physics. It is only stated for the readers
general knowledge.
On the other hand, here is a theorem which is much more useful, due to
Green8:
THEOREM 4.50 (Green-Riemann) Let P and Q be functions from R2 to R or
C, R2 -differentiable, and let K be a sufficiently regular compact subset.9 Then we have
Z
ZZ

Q
P
(P dx + Q dy) =
dx dy.
x
y
K
K
8
George Green (17931841), English mathematician, was born and died in Nottingham,
where he wasnt sheriff, but rather miller (the mill is still functional) and self-taught. He
entered Cambridge at 30, where he got his doctorate four years later. He studied potential
theory, proved the Green-Riemann formula, and developed an English school of mathematical
physics, followed by Thompson (discoverer of the electron), Stokes (see page 472), Rayleigh,
and Maxwell.
9
The boundary of K , denoted K , must be of class C 1 and, for any point z K , there
must exist a neighborhood Vz of z homeomorphic to the unit open ball B(0 ; 1), such that
Vz K corresponds by the homeomorphism with the upper half-plane. This avoids a number
of pathologies.
106
Complex Analysis I
This result can be stated in the language of functions on C in the following

manner.
THEOREM 4.51 (Green) Let F (z, z ) be a function of the variables z and
z , ad-
mitting continuous partial derivatives with respect to each of them, and let K be a
sufficiently regular compact subset of C. Then we have
Z
ZZ
F
F (z, z ) dz = 2i
dA ,
z
K
K
where dA denotes the area integration element in the complex plane. One can also
write F as a function of x and y, which gives us (with a slight abuse of notation)
Z
ZZ
F
F (x, y) dz = 2i
dx dy.
z
K
K
Proof. It is enough to write F = P + i Q andRto apply the formula from Theorem 4.50 to the real and to the imaginary parts of K F dz.
The Green-Riemann formula should be compared with the Stokes formula10 for differential forms, given page 473.
4.3.d Classification of zero sets of holomorphic functions

The rigidity of a holomorphic function is such that, if one knows the values
of f in a very small subset of the open set of definition , then these data
are theoretically sufficient to know f on the whole of (see Exercise 4.1
on page 90). What amount of information actually suffices to determine f
uniquely? The answer is given by the following theorem: a sequence of points
having one accumulation point in is enough.
First we recall the definition of an accumulation point.
DEFINITION 4.52 Let Z be a subset of C. A point c Z is an accumulation
point if there exists a sequence (zn )nN of elements of Z \{c } such that zn c .
(It is important to remark that the definition imposes that zn 6= c for all
n N; without this condition, any point would be an accumulation point, a
constant sequence equal to c converging always to c .)
THEOREM 4.53 (Classification of zero sets) Let be a domain of C (recall that
this means a connected open set) and let f H (). Define the zero set of f by

def
Z( f ) = a ; f (a) = 0 .
10
An approach to complex analysis based on differential forms is presented by Dolbeault

in [29]. It is a more abstract viewpoint at first sight, but it is very rich.
107
If Z( f ) has at least one accumulation point in , then Z( f ) = , that is, the

function is identically zero on .
If Z( f ) has no accumulation point, then for any a Z( f ), there exist a unique
positive integer m(a) and a function g H () which does not vanish at a, such
that one can write f in the factorized form
f (z) = g(z) (z a)m(a)
for all z .
Moreover, Z( f ) is at most countable.

DEFINITION 4.54 The integer m(a) defined in the preceding theorem is called
the order of the zero of f at the point a.
Interpretation
Suppose we have two holomorphic functions f1 and f2 , both defined on the
same domain , and that we know that f1 (z) = f2 (z) for any z belonging
to a certain set Z. Then if Z has at least one accumulation point in , the
functions f1 and f2 are necessarily equal on the whole set .
Thus, if we have a function holomorphic on C and if we know its values for instance on R+ , then it is theoretically entirely determined. In
particular, if f is zero on R+ , then it is zero everywhere.
Note that is not simply open, but it must also be connected. Indeed, if
were the union of two disjoint open subsets, the rigidity of the function
would not be able to bridge the gap between the two open subsets, and the
behavior of f on one of the two open subsets would be entirely uncorrelated
with its behavior on the other.
def
Counterexample 4.55 Assume we know a holomorphic function f at the points zn = 1/n, for
N ,
all n
and that in fact f (zn ) = n. In this case, f is not uniquely determined. In fact,
both functions f1 (z) = 1/z and f2 (z) = exp(2i/z)/z satisfy f1 (zn ) = f2 (zn ) = f (z) for all
n N. What gives? Here, because f (z) when z tends to 0, f is not holomorphic at 0;
the accumulation point, which is precisely 0, does not belong to and the theorem cannot be
applied to conclude.
COROLLARY 4.55.1 Let be a connected open subset and assume that the open ball
B = B(z0 ; r) with r > 0 is contained in . If f H () and if f is zero on B,
then f 0 on .
COROLLARY 4.55.2 Let C be a connected open subset, and let f and g be
holomorphic on and such that f g 0 on . Then either f 0 on or g 0
on .
Proof. Let us suppose that f 6 0; let then z0 be a point of such that f (z0 ) 6= 0.
Since f is continuous, there exists a positive real number r > 0 such that f (z) 6= 0
def
for all points in B = B(z0 ; r ). Hence g 0 on B and Corollary 4.55.1 allows us to
conclude that g 0 on .
Complex Analysis I
108
COROLLARY 4.55.3 Let f and g be entire functions. If f (x) = g(x) for all
x R, then f g.
These theorems also show that functional relations which hold on R remain valid on C. For instance, we have
x R
cos2 x + sin2 x = 1.
The function z 7 cos2 z + sin2 z 1 being holomorpic (it is defined using the
squares of power series with infinite radius of convergence, and is therefore
holomorphic on C) and being identically zero on R, it must be zero on C;
we therefore have
z C
cos2 z + sin2 z = 1.
4.4
Singularities of a function
Certain functions are only holomorphic on an open subset minus one or
a few points, for instance the function z 7 1/z, which belongs to the set
H C \ {0} . It is in general those functions which are useful in physics.
The points at which f is not holomorphic are called singularities and have a
physical significance in many problems.11
4.4.a
Classification of singularities
There are three kinds of singularities: artificial (or removable) singularities,

poles, and essential singularities. The first case is not very interesting. As the
name suggests, it is not a true singularity. It corresponds to the case of a
bounded f . We have the following theorem due to Riemann (see page 138).
THEOREM 4.56 Let a be a point in C. Let f be a function defined on an open
neighborhood of a, except at a, such that f is holomorphic and bounded on \ {a}.
Then lim f (z) exists and the function f can be continued to a function f defined on
by za
putting
def
f (a) = lim f ( ).
a
The function f so defined is holomorphic on the whole of .

11
Thus, in linear response theory, the response function will be, depending on conventions,
analytic (say) in the complex upper half-plane, but will have poles in the lower half-plane.
Those poles correspond to the energies of the various modes. In particle physics, they will be
characteristic of the mass of an excitation (particle) and its half-life.
Singularities of a function
109
DEFINITION 4.57 If a is a singularity of f and f is bounded in a neighbor-
hood of a, then a is called a removable singularity, or an artificial singularity.

Consequently, if f has an artificial singularity, it must have been badly
written to begin with. Consider, for instance, z 7 (4z 2 1)/(2z+1), which has
an artificial singularity at 1/2; we should instead have written the function
as z 7 2z 1, which would then have been defined on the whole of C.
All artificial singularities are not so visible, however, as shown by the next
example.
Example 4.58 The most important example of an artificial singularity is the following. Let f
be a function holomorphic on a domain C. Let z0 and define
def
g(z) =
f (z) f (z0 )
z z0
z \ {z0 }.
Then g is holomorphic on \ {z0 } and has an artificial singularity at z0 . Indeed, if we extend

g by continuity by putting
f (z) f (z0 ) if z 6= z0 ,
def
g(z) =
z z0

f (z0 )
if z = z0 ,
then we obtain a function which is holomorphic on .
Other singularities are classified in two categories.

THEOREM 4.59 (Casorati-Weierstrass)
Let a be an isolated singularity of

a function f H \ {a} . If a is not a removable singularity, then only two
possibilities can occur:

a) lim f (z) exists and is equal to + ;
za
b) for any r > 0 such that B(a ; r) , the image of B(a ; r) \ {a} by f is
dense in C.
Proof. Assume there exists r > 0 such that the image of B(a ; r ) is not dense in C.
Then, for some C and some > 0, we have

| f (z) | >
for all z B(a ; r ) \ {a} .

def
Consider now the function g, defined on B(a ; r )\{a} by g(z) = 1/ f (z) . Then g
is holomorphic on B(a ; r ) \ {a} and, moreover, |g(z)| < 1/ for all z B(a ; r ) \ {a}.
According to the removable singularity theorem, we can continue the function g to a
function g : B(a ; r ) C. This new function g is then nonzero everywhere, except
possibly at a. For all z 6= a, we have
1
f (z) = +
.
()
g (z)
If g (a) were not 0, the right-hand side of () would be holomorphic on B(a ; r ) and
f would have an artificial
singularity at a, which we assumed was not the case. Hence
g (a) = 0 and lim f (z) = +.
za
Complex Analysis I
110
DEFINITION 4.60 (Poles) If f H \ {a} and if

lim f (z) = +,
za
then f admits a pole at a, and a is called a pole of f .

DEFINITION 4.61 (Essential singularity) If f H \ {a} and if f (z)
has no limit as z tends to a, then f is said to have an essential singularity

at a.
Poles and essential singularities are of a very different nature. Whereas
functions behave quite reasonably in the neighborhood of a pole, they are
very wild in the neighborhood of an essential singularity.
Example 4.62 Consider the function f : z 7 exp(1/z). Then, in any neighborhood of 0, the
function f takes any value in C except 0 and oscillates so fast, in fact, that it reaches any
of these complex values infinitely often. This can be seen as follows. Let be an arbitrary
nonzero complex number. There exists a complex number w such that e w = . If we now put
def
zn = 1/(w + 2in), then zn 0 and f (zn ) = exp(1/zn ) = for all n N.
n
4.4.b
Meromorphic functions
DEFINITION 4.63 A subset S of C is locally finite if, for any closed disk D,
the subset D S is finite. If is an open set in C and if S , S is locally
finite in if and only if for any closed disque D contained in , the set
D S is finite.
In practice, to say that S is locally finite in C boils down to saying that

either S is finite, or S is countable and can be written S = {zn ; n N} with
lim |zn | = +.
Example 4.64 Let f be a function holomorphic on an open set and not identically zero.
Then the set of zeros of f is locally finite in .

Example 4.65 The set {1/n ; n N } is locally finite in the open ball B(1 ; 1) (which does not
contain 0), but not in C.
DEFINITION 4.66 Let be an open subset of C and S a subset of . A

function f : \ S C is a meromorphic function on if and only if
i) S is a closed subset of and is locally finite in ;

ii) f is holomorphic on \ S (which is then necessarily closed);
iii) for any z S , f has a pole at z.
Example 4.67 The function z 7 1/z is meromorphic on C. The set of its singularities is the
singleton {0}.
Laurent series
111
PROPOSITION 4.68 If is an open subset of C and if f : C is holomorphic
and not identically zero, then the function

def
F (z) =
1
f (z)
for all z U \ Z( f )
is a meromorphic function.
Proof. The function F is holomorphic on \ Z ( f ) and, since Z ( f ) is locally
finite, according to Theorem

4.53 on page 106, it suffices to check that for any point
a Z ( f ), we have lim F (z) = +, which is immediate by continuity.
za
Example 4.69 If P and Q are coprime polynomials with complex coefficients (i.e., if they have
no common zero in C), with Q 6= 0, the rational function P / Q is meromorphic on C. The

set of its poles is the set of zeros of Q , which is locally finite (in fact finite) because Q is
holomorphic.
Remark 4.70 Two other kinds of singularities are sometimes defined:
Branch points appear when a function cannot be defined on a simply connected open set, as,
for instance, the complex logarithm or other multivalued functions (see Section 5.1.a
on page 135).
Singularities at infinity are defined by making the substitution w = 1/z. One obtains a new
function F : w 7 F (w) = f (1/w); then f has a singularity at infinity if F has a
singularity at 0, those singularities being by definition of the same kind (see page 146).
4.5
Laurent series
4.5.a
Introduction and definition
We have just seen that a function holomorphic on \ z0 could have singularities of two different kinds at z0 : an essential singularity or a pole. We
can shed some light on this distinction in the following manner. If f can be
written in the form of the following series (which is not a power series, since
coefficients with negative indices occur):
X
f (z) =
an (z z0 )n ,
nZ
then
f has a pole at z0 if there are only finitely many nonzero coefficients

with negative indices;
f has an essential singularity at z0 if there are infinitely many.
112
Complex Analysis I
The next theorem stipulates that f usually admits such an expansion.
THEOREM 4.71 (Laurent series) Let f be a function holomorphic on the annulus
1 < |z z0 | < 2 , where 0 1 < 2 +. Then f admits a unique

expansion, in this annulus, of the type
+
X
f (z) =
an (z z0 )n .
n=
Moreover, the coefficients an are given by

Z
1
f ( )
an =
d ,
2i C ( z0 )n+1
where C is any contour inside the annulus turning once around z0 in the positive
direction (i.e., IndC (z0 ) = 1).
Remark 4.72 The annulus considered may also be infinite, as, for instance, C = C \ {z0 }.
This motivates the following definition, due to Laurent12:

DEFINITION 4.73 The preceding expansion is called the Laurent series ex-
pansion of f . If there are only finitely many nonzero coefficients with

negative index in this expansion, then the absolute value of the smallest index with nonzero coefficient is called the order of the pole at z0 . A pole of
order 1 is also called a simple pole.
Proof of Theorem 4.71. Without loss of generality, we may assume that z0 in the
proof. We cut the annulus in the way shown in the figure below. The interior circle,
with radius 1 , is taken with negative orientation, and the exterior circle with positive
orientation. The two parallel segments on the figure are very close, and we will take
the limit where they become identical at the end of the proof. Denote by the
contour thus defined, and by 1 and 2 the two circles with positive orientation. Denote
moreover by the interior of the contour , that is, the open annulus minus a small
strip. This open set is not convex, but it is still simply connected. One can therefore
apply Cauchys formula on (or, more precisely, on a simply connected neighborhood
of containing ).
0
z
We have then, for any point z in the annulus, by Cauchys formula applied to f
and :
Z
Z
Z
1
f ( )
1
f ( )
1
f ( )
f (z) =
d =
d
d
2i z
2i 1 z
2i 2 z
12
Pierre Laurent, (18131853), a former student at the cole Polytechnique, was a hydraulic
engineer and mathematician.
Laurent series
113
in the limit where both segments coincide. It is enough then to write that, for any
1 , we have | z| < 1 and hence
X
1
1
zn
1
=
=
,
n+1
z
1 z/
n=0
whereas for 2 , we have |z/ | < 1 and
1
X
1
1
1
zn
=
=
.
z
z
1 /z
n+1
n=
This proves the first formula stated. Moreover, the value of an is given by the integral
on 1 or 2 depending on whether n is positive or negative. The integral on any
intermediate curve gives of course the same result.
P
an z n converges
for a certain z0 C, then
it
converges
for
any
z
C
with
|z|
<
|z
|.
We
have,
for Laurent
0
P
n
series, a similar result: if +
n= a n z converges for z1 and z2 with |z1 | < |z2 |, then it converges
for any z C such that |z1 | < |z| < |z2 |. Similarly, for a given Laurent series, there exist
two real numbers 1 and 2 (with 2 possibly +) such that the Laurent series converges
uniformly and absolutely in the interior of the open annulus 1 < |z| < 2 and diverges
outside. This result is elementary and does not require the theory of holomorphic functions.
Remark 4.74 We know that a power series has a radius of convergence, that is, if
Remark 4.75 The proof of the theorem also shows that f can be written as f = f1 + f2 , where
f1 is holomorphic in the disc |z| < 2 and f2 is holomorphic in the domain |z| > 1 . If one
asks, in addition, that f1 (z) 0, then this decomposition is unique.
z0
4.5.b
Examples of Laurent series
In the case where a function f has only one singularity, for instance,
1
def
,
f (z) =
z1
the Laurent series may be defined on the (infinite) annulus
def
Indeed, one can write
D1 = {z C ; |z| > 1}.
f (z) =
1
1 1
=
z 1
z 1
1
z
X
1
n
z
n=1
(4.6)
for |z| > 1. Note that this series has infinitely many negative indices, but
beware, this does not mean that there is an essential singularity! Indeed, this
infinite series comes from having performed the expansion in an annulus
around 0, while the pole is at 1. A Laurent series expansion around z = 1
gives of course
1
f (z) =
z1
with no other term. Remark that this Laurent series expansion (4.6) is valid
only in the annulus |z| > 1, outside of which the series diverges rather
trivially.
Complex Analysis I
114
This same function f can be expanded on the annulus 0 < |z| < 1 by
X
f (z) =
zn.
n=0
When a function has more than one singularity, the annulus must remain
between the poles; for example, the function g(z) = 1/(z 1)(z + 2) admits
a Laurent series expansion on the annulus 1 < |z| < 2. Note that it also
admits a (different) Laurent series expansion on the annulus 2 < |z|, and a
third one on 0 < |z| < 1 (this last expansion without any nonzero coefficient
with negative index, since the function g is very cosily holomorphic at 0).
4.5.c The Residue theorem

DEFINITION 4.76 An open subset C is holomorphically simply con-
nected if, for any holomorphic function f : C, there exists a holomorphic function F : C such that F = f .
We will admit the following results:
LEMMA 4.77 An open subset C is holomorphically simply connected if and only
if, for any holomorphic function f : C and for any closed path inside , we
have
Z
f (z) dz = 0.
LEMMA 4.78 An open subset C is holomorphically simply connected if and only
if it is simply connected.
Consequently, we will drop the appellation holomorphically simply connected from now on. The notion of simple connectedness acquires a richer
meaning because of the preceding lemma however.
DEFINITION 4.79 (Residue) Let f be a function that can be expanded in
Laurent series in an annulus 1 < |z z0 | < 2 :
X
f (z) =
an (z z0 )n
for 1 < |z z0 | < 2 .
nZ
The coefficient a1 (corresponding to the term with 1/(z z0 )) is called the

residue of f at z0 and is denoted
a1 = Res ( f ; z0 ).
In particular, if f H ( \ {z0 }) and f has a pole at z0 , we can write
f (z) = g(z) +
N
X
n=1
an
,
(z z0 )n
Laurent series
115
with g holomorphic on , and a1 is the residue of f at z0 .

We will see in the next section how to compute the residue of a function
at a given point.
THEOREM 4.80 Let f be a function holomorphic on the annulus
1 < |z z0 | < 2 .
Let be a closed path contained inside this annulus. Then we have

Z
1
f (z) dz = Ind (z0 ) Res ( f ; z0 ).
2i
(4.7)
Proof. We can indeed write (taking z0 = 0)

X
a
where g(z) =
an z n .
f (z) = 1 + g(z),
z
n6=1
P
P
n
Since the series n6=1 an z converges, as well as the series n6=1 an z n+1 /(n + 1), it
R
follows that the function g has a holomorphic primitive. Hence g(z) dz = 0 and
we get
Z
Z
dz
f (z) dz = a1
= 2i a1 Ind (0),
z
which is the required formula.
Another formulation of this theorem, which will be of more use in practice, is as follows:
THEOREM 4.81 (Residue theorem) Let C be a simply connected
open set
and let P1 , . . . , Pn be distinct points in . Let f H \ {P1 , . . . , Pn } and let

be a closed path contained in which does not contain any of the points P1 , . . . , Pn .
Then we have
Z
n
X
f (z) dz = 2i
Res ( f ; Pi ) Ind (Pi ).
i=1
Why is the Residue theorem useful?

We can make clear immediately that its interest, for a physicist, is mostly
computational. Then it is, in many cases, irreplaceable. It relates a path
integral in the complex plane with a sum which depends only on the behavior
of the function near finitely many points. It is therefore possible to compute
the integral around a closed path (see Section 4.6.d) or, through a small trick,
an integral over R (see Sections 4.6.b and 4.6.c). Conversely, the sum of certain
series can be identified with an infinite sum of residues, which is linked to
an integral over a closed path; by deforming this closed path, one may then
relate this integral to a finite sum of residues (see Section 4.6.e).
But before going into this, we must see how to compute in practice the
residue of a function at a given point.
116
Complex Analysis I
4.5.d Practical computations of residues

In this short section, we will see how to compute in practice the residue of
a function. In almost all cases we will have to consider a function f
H \ {z0 } , where is a neighborhood of a point z0 C. In the case
where f admits a simple pole at z0 , namely, if one can express this function
in the form f (z) = a1 /(z z0 ) + g(z), where g H (), then we have
a1 = Res ( f ; z0 ) = lim (z z0 ) f (z).
zz0
(4.8)
How can one know if f has a simple pole at z0 and not a higher order
pole? If one has no special intuition in the matter, one can try to use the
preceding formula: if it works, this means the pole was simple; if it explodes,
the pole was of higher order!13
Consider, for example, the case where f is a function that has a pole of
order 2 at z0 . Then we have
a2
a1
f (z) =
+
+ g(z)
with g H ().
(4.9)
2
(z z0 )
(z z0 )
Formula (4.8) then gives indeed an infinite value in this situation. How, then,
can one obtain a1 without being annoyed by the more divergent part? If
we multiply equation (4.9) by (z z0 )2 , we get
(z z0 )2 f (z) = a2 + a1 (z z0 ) + (z z0 )2 g(z).
Differentiate with respect to z:

d
(z z0 )2 f (z) = a1 + 2(z z0 ) g(z) + (z z0 )2 g (z).
dz
By taking the limit [z z0 ], we obtain exactly a1 and this without explosion. As a general rule, for a pole of order k, one can use the following
formula, the proof of which is immediate:

THEOREM 4.82 Let f H \ {z0 } be such that f has a pole of order k at z0 .
Then the residue of f at z0 is given by
o
1
dk1 n
Res ( f ; z0 ) =
lim k1 (z z0 )k f (z) .
(4.10)
(k 1)! zz0 dz
Example 4.83 Let a be a nonzero complex number. Consider the function f (z) = 1/(z 2 + a 2 ),
which is meromorphic on the complex plane and has two poles at the points +ia and ia.
These poles are simple and we have
(z ia)
Res ( f ; ia) = lim (z ia) f (z) = lim
zi a
zi a (z ia)(z + ia)
1
1
= lim
=
.
zi a (z + ia)
2ia
13
I agree that this is a dangerous method, but as long as it is restricted to mathematics and
is not applied to chemistry, things will be fine.
Application to the computation of horrifying integrals. . .
117
4.6
Applications to the computation of horrifying
integrals or ghastly sums
The goal of this section is to show how the knowledge of Cauchys theorems and residues can help computing some integrals which are otherwise
quite difficult to evaluate explicitly. In certain cases one can also evaluate the
sum of a series using similar techniques.
The idea, when trying to compute an integral on R, is to start with only
a piece of the integral, on an interval [R, R], then close this path in the
complex plane, most often by a semicircle, checking that the added piece does
not contribute to the integral (at least in a suitable limit), and compute this
new path integral by the Residue theorem. Before presenting examples, we
must establish conditions under which the extra piece of the integral will
indeed be negligible.
4.6.a
Jordans lemmas
Jordans lemmas are little theorems which are constantly useful when trying
to compute integrals using the method of residues.
The first of Jordans lemmas is useful for the computation of arbitrary
integrals.
THEOREM 4.84 (First Jordan lemma) Let f : C C be a continuous function
in the sector

def
S = r e i ; r > 0 and 0 1 2 ,
such that lim z f (z) = 0 (resp. lim z f (z) = 0).

|z|0
|z|

Denote (r) = r e i ; 1 2 . Then we have
Z
Z
lim+
f (z) dz = 0
(resp. lim
f (z) dz = 0).
r 0
r +
(r )
(r )
2
1
Complex Analysis I
118
The second lemma is more specifically useful to compute Fourier transforms.

THEOREM 4.85 (Second Jordan lemma) Let f : C C be a continuous func-
tion on the sector

S = r e i ; r 0 and 0 1 2 ,
such that z
lim f (z) = 0. Then
zS
lim
r +
f (z) e i z dz = 0.
(r )
Proof. Consider the case (the most difficult) where the sector is 0 . Remarking that

i re i
e
= e r sin ,
we can bound the integral under consideration by
Z
Z

iz
f (r e i ) r e r sin d
f
(z)
e
dz

(r)

0
Z /2

f (r e i ) + f (r e i () ) r e r sin d.
Let > 0. Since f is a function that tends to 0 at infinity,

there exists R > 0 such
that, for any r R and for any [0, ], we have f (r e i ) . In particular, for
any t R, we have
Z

Z /2

iz
f (z) e dz 2
r e r sin d.

(r)

0
But for any 0, 2 , we have 0 2

sin . We deduce that
Z /2
Z /2
r e r sin
r e 2r/ d = (1 e r ) ,
2
2
0
0
which shows that for any r R, we have
Z

iz
f (z) e dz .

(r)

4.6.b Integrals on R of a rational function

We consider the case where F is a rational function (i.e., a function of the
form F (x) = P (x)/ Q (x), where P and Q are two polynomials,
which we
R
assume have no common zeros), and we wish to evaluate R F (x) dx. We will
assume here that the function F has no real pole, i.e., the polynomial Q does
not vanish on R (the situation where Q vanishes will be treated in Chapter 8,
during the study of Cauchys principal values, page 223).
119
The nephew of the symbolist painter Puvis de Chavanne, Camille

Jordan (18381922), born in Lyons, was first in the entrance exam
for the cole Polytechnique at 17, then entered the cole des
Mines and became engineer in charge of Parisian quarries. He
taught at the cole Polytechnique, then entered Collge de France
and the Acadmie des sciences. He studied linear algebra, group
theory (he introduced the terminology abelian groups and studied cristallography; among his students were Felix Klein and Sophus Lie), geometry, and Fourier analysis, and made some of the
first steps in measure theory. He introduced some fundamental
topological concepts (homotopy). His pedagogical influence was
enormous.
We will also assume that the integral in question is indeed convergent,

which is the case if the degrees of P and Q satisfy deg Q deg P + 2.
We start by extending the domain of the real variable to the complex plane,
then truncate the integral on R to an integral on [R, R], with R > 0, and
close the path of integration by means of a semicircle CR of radius R as in
Figure 4.4.
Denote by (R) = [R, R] CR the contour thus defined. According
to the assumptions,
R we can use the first of Jordans lemmas 4.84, with the
consequence that C F (z) dz 0 and hence
R
lim
R+ (R)
F (z) dz =
F (x) dx.
Moreover, when R goes to infinity, the contour (R) ends up containing all
the poles of F in the upper half-plane, i.e., for R large enough, we have
Z
X
F (z) dz = 2i
Res (F ; a),
(R)
a pole in
upper half-plane
which gives, in the limit R +, the desired integral

Z
F (x) dx = 2i
Res (F ; a).
(4.11)
a pole in
upper half-plane
Remark 4.86 We might, of course, have used a different contour in the argument to close the
segment [R, R] in the lower half-plane. Then we would have had to compute the residues
in the lower half-plane. But then one would have had to be careful that with the orientation
chosen the winding number of the contour around each pole would be equal to 1 and not
1, since the contour would have been oriented negatively.
Complex Analysis I
120
Fig. 4.4 Integration path for a rational function.
4.6.c
Fourier integrals
We are looking here at evaluating integrals of the type

Z +
I =
f (x) e ik x dx
with k R,
which occur very often in practice, since they correspond to the value of the
Fourier transform of f evaluated at k. One can use the same type of argument
and contour as the one used in Section 4.6.b, using this time the second of
Jordans lemmas, but one must be careful of the sign of k! Indeed, the integral
on the upper semicircle

def
C + = R e i ; [0, ] =
(4.12)
R
of f (z) e ikz only tends to 0 when r tends to infinity when the real part of ikz
is negative; if z C + , we must then have k positive. If k is negative, we must
rather perform the integration along the lower semicircle

C = R e i ; [0, ] =
def
Example 4.87 We try to compute

def
J =
(4.13)
cos(k x)
dx,
x 2 + a2
where a is a strictly positive real number. We start by writing the cosine in exponential form.
One way to do it is to use cos(k x) = (e i kx + e i kx )/2, but here it is simpler to notice that
Z
e i kx
J =
dx,
2
2
x + a
the imaginary part being zero. For this last integral, we must now distinguish the two cases
k > 0 and k < 0.
121
When k > 0, we take the contour (4.12) and we obtain

i kz

Z
e
e i kz
dz
=
2i
Res
;
ia
,
J = lim
2
2
R
z 2 + a2
+ (R) z + a
since the only poles of the function being integrated are ia and ia. These poles are
simple (because 1/(z 2 + a 2 ) = 1/(z + ia)(z ia)), and we compute easily
i kz

e
e i kz
e i kz
e ka
Res
;
ia
= lim (z ia) 2
= lim
=
,
2
2
2
zi
a
zi
a
z +a
z +a
z + ia
2ia
which gives
ka
e .
a
When k < 0, we must take the contour (4.13) and we get
Z
i kz

e i kz
e
J = lim
dz
=
2i
Res
;
ia
.
2
2
R
z 2 + a2
(R) z + a
J =
We then evaluate the residue at ia, to obtain
J = e ka .
a
To conclude, we have, independently of the sign of k (and even if k = 0)
Z
cos k x
dx = e |k|a .
2 + a2
x
a
4.6.d
Integral on the unit circle of a rational function
We are trying here to compute an integral of the type

Z 2
K =
R(sin t, cos t) dt,
0
where R(x, y) is a rational function which has no pole on the unit circle
(x 2 + y 2 = 1). Putting, naturally, z = e i t , we obtain the relations
e i t e i t
1
1
1
1
=
z
, cos t =
z+
.
dz = ie i t dt, sin t =
2i
2i
z
2
z
The integral becomes
K =
R
C
z 1/z z + 1/z
,
2i
2
dz
,
iz
where C is the unit circle with the positive orientation. An immediate application of the Residue theorem shows that
K = 2i
a pole B(0 ; 1)
Res

1
1
1 1
1
R
z , z+
;a
iz
2i
z 2
z
(4.14)
Complex Analysis I
122
Example 4.88 We want to compute the following integral:
I=
dt
,
a + sin t
a R, a > 1.
Put R(x, y) = 1/(a + x). Then we have

I = 2i
1
iz a +
Res
a ple B(0 ; 1)
def
1
2i
1
z
; a .
Define now f (z) = 2/(z 2 + 2iaz 1), and it only remains to find the poles of f . Thesepare the
solutions of the equation
z 2 + 2iaz 1 = 0, and so they are the points = ia + i a 2 1
p
and = ia i a 2 1. They are located as shown in the figure below
p
2
Only the pole contributes to the integral,
p it is a simple pole with residue i/ a 1
and the integral is therefore equal to I = 2/ a 2 1.
4.6.e Computation of infinite sums

Consider a series of the form
S=
f (n),
nZ
where f is a given meromorphic function, decaying sufficiently rapidly for

the sum to exist, and with no pole on the real axis. One can check easily that

f (n) = Res f (z) cotan z ; z = n ,
and so
S=
X
nZ
f (n) =
nZ

Res f (z) cotan(z) ; z = n .
The partial sum from N to N is therefore given by the integral of f along

the following contour14 (where each circle is taken with positive orientation):
which we can deform (adding contributions which cancel each other in the
integral) into
14
Which is not a contour properly speaking since it consists of multiple pieces, but this is
without importance, the integral on this contour being defined as the sum of the integrals
on each circle.
123
and by continuity into
In the limit where N tends to infinity, and if the function f decays to 0, the
contour becomes simply
Finally, this last integral can be computed by the method of residues, by using
Jordans Lemmas and closing each of the lines by semicircles. There only
remains to evaluate the residues of the function f outside the real axis.
Example 4.89 Suppose we want to evaluate the sum
S=
X
nZ
1
,
n2 + a 2
where a is a nonzero real number. We will then define

1
def
f (z) = 2
cotan(z),
z + a2
and S is the sum of the residues of this function at all the integers if a 6= 0. We will therefore
assume a > 0. The reasoning explained above shows that this sum is equal to the opposite of
the sum of the residues of f at nonreal poles. These are located at z = ia, and we have
Res ( f ; z = ia) =
1
cotan(ia) = coth(a),
2ia
2a
hence the result
coth(a).
(4.15)
a
2
One can note that, if a tends to 0, we have S 1/a , which is the expected behavior (the term
n = 0 of the sum is 1/a 2 ). In the same manner, we have
S=
X
n=1
1
=
coth(a) 2 .
n2 + a 2
2a
2a
The method must obviously be tailored to each case, as the next example
shows:
Example 4.90 We now try to evaluate the sum
T =
X
(1)n
n=1
n4
Instead of using the function cotan, it is interesting to use the function z 7 1/ sin z, which
has poles at every integer and for which the residue at n N is precisely equal to (1)n /.
Complex Analysis I
124
Unfortunately, the function z 7 1/z 4 has a pole at 0. This pole must then be treated separately.
We start by noticing that
T =
T
2
T =
with
X
(1)n
,
4
n= n
n6=0
and we deduce that the sum T is given by 1/2i times the integral on the following contour
of the function z 7 1/z 4 sin(z)
which we then close (appealing to Jordans second lemma) into
before deforming it by continuity into
which gives
T = Res
1
z 4 sin z
; z=0 ,
the minus sign coming from the fact that the last path is taken with negative orientation. There
only remains to evaluate this residue, using, for instance, a Taylor expansion of 1/ sin z:
z2
z5
sin z = z
+
+ ,
6
120

1
1
2 z 2 4 z 4
1
=
1
+
+
z 4 sin z
z5
6
120

2
2
1
z
144 z 4
= 5 1+
+
+ ,
z
6
720
which shows that the residue of 1/z 4 sin z at 0 is
X
(1)n
n=1
n4
144
720
and that
74
.
720
Exercises
125
EXERCISES
Exercise 4.4 Show that if f is holomorphic on an open subset C, then
2
f

| f |2 = 4 .
z
Exercise 4.5 Let f and g be continuously differentiable functions on R2 identified with C
(i.e., functions such that the partial derivatives f / x and f / y exist and are continuous).
Show that the following chain rules hold:
f g
f g
( f g) =
+
,
z
z z
z z
f g
f g
( f g) =
+
.
z
z z
z z
and
Exercise 4.6 Let U be a connected open set and let f : U C be a function such that f 2
and f 3 are holomorphic. Show that f is holomorphic on U .
Path integration
Exercise 4.7 Consider, in C, the segments T1 = [i, 6i], T2 = [i, 2 + 6i] and T3 =
2
[2
R + 5i, 6i]. Define
R now a function f : C C by f (x + i y) = x 3i y. Compute explicitly
f
(z)
dz
and
f
(z)
dz.
Conclude.
T
T T
1
Exercise 4.8 Denoting by T the circular arc joining 3 to i
compute
p
3 with equation
(x 1)2 + y 2 = 4,
dz/z 2 .
Exercise 4.9 Recall that
a R, the integral
R +
e x dx =
Z
. One then wishes to compute, for any nonzero

2
e x cos(a x) dx.
Let f (z) = e z for any z C, and denote by T the rectangle

Compute the integral
T = [R, R, R + ia/2, R + ia/2].
f (z) dz and conclude by letting R go to infinity.
Integrals using residues

Exercise 4.10 Compute the residue at 0 of sin(z)/z 2 .
Exercise 4.11 Prove the convergence of the following integrals and compute their values:
(1 + x) sin 2x
dx
x 2 + 2x + 2
and
dx
(1 + x 2 )n
with n N
Complex Analysis I
126
Exercise 4.12 (Fresnel integral) Does the Fresnel integral

def
I =
e i x dx
exist in the sense of Lebesgue?

We nevertheless try to assign a meaning to this integral. Show that
Z R
def
2
I (R, R ) =
e i x dx
R
exists and has a limit as R and

go to +, and compute this limit, using a suitable contour
and remembering the following value of the gaussian integral:
Z +
p
2
e x dx = .
Exercise 4.13 (Fourier transform of the gaussian) Let R. Compute
e x e 2i x dx,
which is the Fourier transform of the gaussian x 7 e x at the point . One may use
Exercise 4.9.
(Another method will be given in Exercise 10.1 on page 295, using the Fourier transform.)
Exercise 4.14 Let A be a finite subset of C containing no positive real number. Let
C \ Z. Let finally f : C C be a function meromorphic on C and holomorphic on C \ A,

such that
Z 2

def
f (r e i ) r Re+1 d.
lim Ir = 0
with
Ir =
r0
or r+
Show that we have then

Z +
f (x) x dx =
0
2i X
Res
1 e 2i aA

f (z) e L(z) ; a ,
where L denotes the complex logarithm defined on C \ R+ by

def
L(z) = log |z| + i(z)
with (z) ]0, 2[ .
Exercise 4.15 Let f be a function holomorphic in the open disc D(0; R) with R > 1.

def
Denote by the unit circle: = e i ; [0, 2[ .
Computing in two different ways the integrals
Z
Z
f (z)
f (z)
I=
2 + z + 1/ z
dz,
J =
2 z + 1/ z
dz,
z
z
show that

2 2
f e i cos2 /2 d = 2 f (0) + f (0),
Z 2

2
f e i sin2 /2 d = 2 f (0) f (0).
0
Exercises
Exercise 4.16 Compute the value of
127
dx
.
x6 + 1
On may use the classical contour or, more cleverly, a sector with angle /3.
Exercise 4.17 Compute
log x
dx.
1 + x3
Exercise 4.18 Compute
cos
x
2
x2 1
+
2
sin x dx =
dx.
cos x 2 dx =
1
2
.
2
Hint: Use the following contour: from 0 to R on the real axis, then an eighth of a circle
up to R e i /4 , and then back to the starting point along the imaginary axis.
+
0
cosh a x
,
dx =
cosh x
2 cos 2 a
where |a| < 1.
e az
, and consider the contour integral on the
cosh z
rectangle with vertices at [R, R, R + i, R + i].
Hint: Look for the poles of F (z) =
Exercise 4.21 Show that if R is a rational function without poles on the half-axis of positive
real numbers R+ , if ]0, 1[ and if R(x)/x is integrable on R+ , then

Z +

X

R(x)
R(z)
1 e 2i
dx
=
2i
Res
,
x
z
0
poles other
than z = 0
where the function z 7
is defined by
ei
e i
for all [0, 2[.
Hint: Use the following contour and take

the limits [R +] and [ 0].
Exercise 4.22 Show that if ]0, 1[, then
dx
=
.
x (1 + x)
sin
Complex Analysis I
128
Exercise 4.23 Let a be a positive real number. Compute then
lim+
dx
.
(x 2 a 2 i)
Hint: Take the following contour and the

limit [R +].
Sums using residues

Exercise 4.24 Compute the following two series:
S=
nN
1
n4 + a 4
and
T =
nN
n2
n4 + a 4
with a R, a 0.
Exercise 4.25 (Ramanujans formula) Using the function z 7 cotan z, deduce, by means
of a clever contour, that
coth coth 2
coth n
197
+
+ +
+ =
.
7
7
7
1
2
n
56 700
Note: This beautiful formula is due to the great indian mathematician Ramanujan.15 But
Ramanujan did not use the calculus of residues.
PROBLEM
Problem 2 Application to finite temperature field theory In quantum statistical me-
chanics, some quantities at equilibrium (such as the density of a gas, its pressure, its average
energy) are computed using sums over discrete frequencies (called Matsubara frequencies).
We consider here a simple case where these sums are easily computable using the method of
residues.
For instance, if one considers a gas of free (i.e., noninteracting) electrons (and positrons16),
denote by the chemical potential of the electrons; the positrons have chemical potential
equal to . The temperature being equal to T , put = 1/kb T , where kb is the Boltzmann
constant.
The free charge density, in momentum space, is given [55, 60] by
X
4(il + )
def
,
N ( p) =
2
2 2
(i
l + ) E p
l
15
Srinivasa Ramanujan (18871920) was probably the most romantic character in the mathematical pantheon. An Indian mathematician of genius, but lacking in formal instruction, he
was discovered by G. H. Hardy, to whom he had mailed page after page covered with formulas
each more incredible than the last. The reader may read Ramanujan, an Indian Mathematician
by Hardy, in the book [46]. One will also find there the famous anecdote of the number
1729.
16
In relativistic quantum mechanics, one cannot have a description of matter which is purely
electronic, for instance; positrons necessarily occur, as was discovered by Dirac [49].
with
l =
(2l + 1)
129
def
E p = (m 2 c 4 + p 2 c 2 )1/2 .
and
1. Put this sum in the form of a sum of residues of a well-chosen function f (z). Make a
drawing of the complex plane with the corresponding poles.
2. Does the chosen function admit any other poles?
3. By considering the decay of the function f (z) at infinity, show that the sum of the
residues vanishes.
4. Show then that we have
N ( p) = tanh
E p
2
+ tanh
= 2N Fel ( p) 2N Fpos ( p),
E p +
2
where we use the electronic and positronic Fermi distributions

N Fel ( p) =
1
e (E p ) + 1
N Fpos ( p) =
and
1
.
e (E p +) + 1
Interpret these two terms.
SOLUTIONS
Solution of exercise 4.1 on page 90. Put f = P + i Q .
We start by remarking that property (ii) implies all the others. Let us show the equivalence
of (ii) and (iii). Suppose, for example, that P = Cnt . Then, at any point,
P
P
=
= 0,
x
y
which, by the Cauchy equations, implies that
Q
Q
=
=0
y
x
and hence that Q = Cnt . One deduces that (ii) and (iii) are equivalent. Moreover, (ii) and (iii),
together, imply property (i). Thus we have shown the equivalence of (i), (ii), and (iii).
Let us show that (iv) implies (iii) this will prove also that (iv) implies (i) and (ii). Assume
then property (iv). We can also assume that f is not the zero function (which is an obvious
case). Then f has no zero at all, and so the sum P 2 + Q 2 = Cnt is never zero. We have
0=
| f |2
P
Q
= 2P
+ 2Q
x
x
x
Cauchy
= 2P
Q
Q
+ 2Q
y
x
and
| f |2
P
Q
Q
Q
= 2P
+ 2Q
= 2P
+ 2Q
.
y
y
y
x
y
The two preceeding equations can therefore be written in the form
Q
P
Q / x
= 0.
P Q
Q / y
0=
But the determinant of this matrix is P 2 + Q 2 which is everywhere non-zero. So we deduce

that, at any point, we have Qx = Qy = 0 and, consequently, that Q = Cnt .
Complex Analysis I
130
There only remains to show that (v) implies (ii). Assume that both f and f are holomorphic;
then the Cauchy equations allow us to write
( f H ) =
Q
P
=
x
y
( f H ) =
and
P
( Q )
=
,
x
y
and hence P / x = 0; similarly P / y = 0, so that P = Cnt . Which, taking into account

the remarks above, completes the proof.
Solution of exercise 4.4. One

checks first that the laplacian can be written = 4. Then
write | f |2 = 4( f f ) = 4 ( f ) f + f ( f ) . Since f is holomorphic, we have f = 0

hence | f |2 = 4 f f = 4 f f since f is an antiholomorphic function; thus f = 0.
To conclude, one finds
| f |2 = 4 f f = 4 | f |2 .
Solution of exercise 4.6. Disregarding the trivial case where f = 0, write then f = f 3 / f 2
at any point where f does not vanish. But since f 2 is holomorphic, its zeros are isolated, and
therefore those of f also, since they are the same. Denote by Z the set of these zeros. Then
f is holomorphic on \ Z and bounded in the neighborhood of each point in Z , which are
therefore artificial singularities. Hence, if we extend by continuity the function g = f 3 / f 2 (by
putting g(z) = 0 for any x Z ), the resulting function g is holomorphic by Theorem 4.56 on
page 108, and moreover g = f by construction.
2
Solution
of exercise 4.8. The

function z 7 (1/z ) being holomorphic on the open set
= z C ; Im(z) > Re(z) which contains the contour T , it admits a primitive on ,

which is evidently z 7 1/z. Then
Z
p
i
1
1 i 3 1
dz
=
= +p .
2
z
z
3
3
3
T
Solution of exercise 4.9. Show that the integrals on the vertical segments tend to 0 when
[R ]. Show moreover that
ZR
Z
2
2
e (x+i a/2) dx = e a
R
e x cos a x dx
by parity. The integral considered is equal to
p a2
e .
Solution of exercise 4.10. The Laurent series expansion of the sine function coincides with
its Taylor series expansion and is given by
z3
z5
+
;
hence
3!
5!
and the residue is the coefficient of 1/z, namely 1.
sin z = z
sin z
1
z
z3
= +
,
2
z
z
3!
5!
Solution of exercise 4.11. For the first integral, put

f (z) =
(1 + z) e 2i z
z 2 + 2z + 2
and integrate by closing the contour from above (with justification!). The only residue
inside
R
the contour is the one in z = 1 + i. The desired integral is the imaginary part of f (z) dz =
ie 22i , which is ( cos 2)/e 2 .
For the second integral, close the contour from above, using Jordans second lemma, with
f (z) =
1
1
=
.
(1 + z 2 )n
(z i)2 (z + i)n
131
Then the desired integral is equal to 2i Res ( f ; i) (or, if the contour is closed from below, to
2i Res ( f ; i)). But, according to formula (4.10), page 116, since i is a pole of order n,

dn1
1
1
.
Res ( f ; i) =
(n 1)! dz n1 (z + i)n z=i
An easy induction shows that
dn1
1
(n) (n 1) (2n + 2)
(1)n1 (2n 2)!
=
=
,
dz n1 (z + i)n
(z + i)2n1
(z + i)2n1 (n 1)!
which leads, after short algebraic manipulations, to

Z +
dz
(2n 2)!
= 2n2
2 .
2 n
2
(n 1)!
(1 + z )
4.12. The Fresnel integral is not defined in the Lebesgue sense since
2
the modulus of e i x is constant and nonzero, and therefore not integrable. However, the
integral is defined as an improper integral, which means that the integral I (R, R ) admits a
limit as R and R (separately) tend to infinity. Consider now the contour
R
R
def
C (R, R ) =
R
The integral on the upper eighth arc R is, by a change of variable, equal to
Z
Z
d
2
e i z dz =
ei p ,
R
C
where C is the upper quarter-circle with radius R 2 , and this last integral does tend to 0 by
Jordans second lemma. The second eighth of the circle is treated in the same manner. Since
the integral on C (R, R ) vanishes (the function is holomorphic), the Fresnel integral is equal
to an integral on the line segment with angle /4:
Z R
Z R
i x2
i (e i /4 x)2 i /4
lim
e
dx
e
e
dx
= 0,
R,R +
hence
e i x dx = e i /4
1+ip
2
.
e x dx = p
2
Solution of exercise 4.13. Denote by F () the integral under consideration. It can be put
in the form
F () = e
e (x+i ) dx.
Denote by the line with imaginary part , hence parallel to the real axis. Then, using
2
Cauchys theorem for the function z 7 e z , which is holomorphic on C, show as in
exercise 4.9, that
Z
Z
2
e z dz =
which gives F () = e .
e x dx = 1,
Complex Analysis I
132
Solution
of exercise 4.16. Consider the contour (R) made from the line segments [0, R],
0, R e i /3 and the part of the circle C (0 ; R) with argument varying from 0 and /3. Let
g be the function g : z 7 z 6 + 1. Then the function f = 1/g has only one pole inside this
contour, at the point = exp(i/6). Check then (using the first Jordan lemma and a few
algberaic computations) that
Z
Z
+ dx
.
f (z) dz 1 e i /3
R+
x6 + 1
(R)
0
Using the residue theorem, it follows therefore that

Z +
dx
2i Res ( f ; )
=
.
6+1
x
1 e i /3
0
To compute the preceeding residue, one may notice that
Res ( f ; ) = lim
z
z
z
= lim
z g(z) g( )
g(z)
1
e 5/6
=
=
.
g ( )
6
Z
Thus we get
because g( ) = 6 1 = 0
i e 5/6
dx
=
= .
+1
3 1 e i /3
3
x6
Solution of exercise 4.17. Define

f (z) =
log z
,
1 + z3
where the logarithm function is defined on C \ R , for instance. Integrate on the contour
going from to R (for > 0), followed by an arc with angle 2/3, then by a line segment
toward the origin, but finished by another arc of the circle with radius . Show that the
integrals on the circles tend to 0 when R tends to infinity and to 0. The only residue inside
the contour is from the pole at z = e i /3 and is equal to 9 e i /6 . Compute the integral on
the second side in terms of the desired integral and in terms of
Z +
2
1
dx = p
1 + x3
3 3
0
(which can be computed using the partial fraction expansion of 1/(1 + x 3 ), for instance).
One finds then
Z +
log x
2
dx = 2 .
3
1
+
x
27
0
Solution of exercise 4.18. Show that the function is integrable using a Taylor expansion at
1 and 1. Integrate on a contour which consists of the line segment from R to R, avoiding
the poles at 1 and 1 by half-circles above the real axis with small radius , and close the
contour by a half-circle of radius R. Show that the integral on the large circle of
f (z) =
e i z/2
z2 1
tends to 0 and compute the integrals on the small circles (which do not vanish; their sum is
equal to in the limit where [ 0]).
Z +
Hence one finds that
f (x) dx = and the desired integral has value /2.
133
Solution of exercise 4.24. Notice that

T + ia 2 S =
X
n=1
1
.
(n2 ia 2 )
Use then the result given by equation (4.15), page 123, and take the real and imaginary parts,
respectively.
Solution of exercise 4.25. Let

f (z) =
cotan(z) coth(z)
.
z7
This function has poles at n and in, for any n Z. The residue at n 6= 0 is simply coth(n)/n7 ,
and the residue at in 6= 0 is cotan(in)/(in)7 = coth(n)/n7 , which gives us that the integral
of f on the square with center at 0 and edges parallel to the axis at distance R = (n + 21 )
(which does not pass through any pole) is
In = 2i
k=n
X
2 coth k
+ 2i Res ( f ; 0).
k7
k=n
k6=0
Since (by an adaptation to this situation of Jordans second lemma) In tends to 0 when n tends
to infinity, we deduce that
S=
n
X
coth k
k=1
k7
1
= Res ( f ; 0).
4
It suffices now to expand coth z cotan z to order 6 around 0 (by hand for the most
courageous, with a computer for the others):
coth z cotan z =
1
73 2
197 6
z
z + O(z 10 )
2
z
45
14 175
which indicates that the residue of f at 0 is 197 /14 175, and the formula stated follows.
Complex Analysis I
134
Sidebar 3 (Differentiability of a function on R 2 ) Consider

f :
R2 R2 ,
(x, y) 7 f (x, y),
a function of two real variables with values in R2 . Denote its coordinates by f = ( f1 , f2 );

thus the functions f1 and f2 take values in R.
The function f is differentiable at a point (x0 , y0 ) R2 if there exists a linear map
L (R2 ) which satisfies
f (x0 + h, y0 + k) = f (x0 , y0 ) + (h, k) + o(h, k)
[(h, k) (0, 0)].
Such a linear map is called the differential of f at the point ( x0 , y0 ), and is

denoted d f(x0 , y0 ) . It can be represented as a 2 2 matrix:
fx fx
x
y
,
mat d f(x0 , y0 ) =
f
y fy
x
y
which means that
f x (x0 + h, y0 + k)
f y (x0 + h, y0 + k)
fx fx
!
f x (x0 , y0 )
x
y h
+
+ o(h, k)
f
f y (x0 , y0 )
y fy k
x
y
fx
fx
!
h+
k
f x (x0 , y0 )
x
y
+ o(h, k).
+
f
f y (x0 , y0 )
y h + f y k
x
y
!
The map which associates d f(x0 , y0 ) to (x0 , y0 ) is called the differential of f and is denoted d f .
It is customary to write dx : (h, k) 7 h and dy : (h, k) 7 k, and consequently the
differential d f can also be written, in the complex notation f = f x + i f y , as
df =
f
f
dx +
dy.
x
y
()
(Notice that we began with a vector function ( f x , f y ) and ended with a scalar function
f x + i f y ; therefore the matrix of d f is a simple line 1 2 and no more a 2 2 square
matrix. It is now the matrix of a linear form.)
If f is differentiable, then its partial derivatives are well-defined. The converse is not true,
but still, it the partial derivatives of f are defined and continuous on an open set R2 ,
then f is indeed differentiable, and its differential d f satisfies the equation ().
Chapter
Complex Analysis II
5.1
Complex logarithm; multivalued functions
5.1.a
The complex logarithms
We seek to extend the real logarithm function x 7 log x, defined on R+ , to

the complex plane or at least to a part, as large as possible, of the complex
plane.
The most natural idea is to come back to one of the possible definitions
of the real logarithm: it is the primitive of x 7 1/x which vanishes at x = 1.
Let us therefore try to primitivize the function z 7 1/z and so put, for any
z C ,
Z
1
def
L(z) =
d ,
(5.1)
(z)
where (z) is a path in the complex plane joining 1 to z, not going through 0
(so that the integral exists). Does equation (5.1) really define a function? The
answer is yes if any choice of the path leads to the same value of L(z),
and no otherwise. However, two paths 1 and 2 , joining 1 to z, but going
on opposite sides of the point 0, give different results (see Figure 5.1); if we
call the path consisting of 1 followed by 2 , the integral of 1/ around
is equal to 2i times the residue
R at 0 (which is 1), times the number of turns
of around 0; so, in the end, (1/ ) d = 2i.
Complex Analysis II
136
2
Fig. 5.1 It is not possible to integrate the function z 7 1/z on the complex plane, since
the result depends on the chosen path. Thus, the integrals on the solid contour
1 and the dashed contour 2 differ by 2i.
Consequently, it is not possible to define the complex logarithm on the

punctured plane C = C \ {0}.
Still, we may restrict our ambitions a bit and only consider the plane from
which a half-line has been removed, namely C \ R .
PROPOSITION 5.1 The complex plane minus the half-line of negative real numbers
C \ R is simply connected.
A consequence of this theorem is that, on this simply connected set, all

closed paths are homotopic1 to a point, and the relation (5.1) does define a
function correctly.
DEFINITION 5.2 The function L(z) defined by
def
L(z) =
1
d ,
(z)
(5.2)
where (z) is a path inside C \ R joining 1 to z, is called the principal

complex logarithm.
Notice that one could just as well remove another half-line from the complex plane, for instance, the half-line of positive real numbers, or one of the
half-lines of the purely imaginary axis. Even more, it would be possible to take
an arbitrary simple curve C joining 0 to infinity; then, on the complement
C \ C , which is simply connected, one can define a new logarithm function.
DEFINITION 5.3 A cut is any simple curve joining 0 to infinity in the com-
plex plane. To each cut C is associated a logarithm function C on C \ C .

The point 0 is called a branch point.
The various logarithm functions are of course related to each other by the
following result:
1
I.e., a closed path can be contracted continuously to a point, see page 575.
Complex logarithm; multivalued functions
137
THEOREM 5.4 Let and be two complex logarithms, defined on C \ C and
C \ C , respectively. Then, for any z C \ {C C }, the difference (z) (z) is

an integral multiple of 2i.
This difference is constant on each connected component of C \ {C C }.
In addition, the principal logarithm satisfies the following property:
PROPOSITION 5.5 If z C is written z = e i , with ], [ and
R+ , then L(z) = log + i.
Example 5.6 Let (z) be the complex logarithm defined by the cut R+ . The reader will easily
check that, if z = e i with ]0, 2[ this time and R+ , then (z) = log + i.
5.1.b
The square root function
The logarithm is not the only function that poses problems when one tries to
extend it from the (half) real line to the complex plane. The square root2
function is another.
Indeed, if z = e i , one may be tempted to define the square root of z by
p def p i/2
z = e . But a problem arises: the argument is only defined up to
2, and so /2pis only defined up to a multiple of ; thus, there is ambiguity
on the sign of z.
It is therefore necessary to specify which determination of the argument
will be used to define the square root function. To this end, we introduce
again a cut and choose a continuous determination of the argument (which
always exists on a simply connected open set).
In the following, each time the complex square root function is mentioned,
it will be necessary to indicate, first which is the cut, and second, which is the
determination of the argument considered.
Example 5.7 If the principal determination of the argument = Arg z ], [ is chosen,
p def p
one can defined z 7 z = e i /2 and the restriction of this square root function to the
real axis is the usual real square root.
By chosing iR as the cut, one can define a square root function on the whole real axis.
5.1.c Multivalued functions, Riemann surfaces

There is a trick to avoid having to introduce a cut for functions like the logarithm or the square root. Physicists call these in general multivalued functions,
and mathematicians speak of functions defined on a Riemann surface.
We start with this second point of view.
The main idea is to replace the punctured complex plane (i.e., without
0), on which the definition of a holomorphic logarithm or square root was
p
The notation to designate the square root (this symbol is a deformation of the
letter r) was introduced by the German mathematician Christoff Rudolff (15001545).
2
Complex Analysis II
138
Georg Friedrich Bernhard Riemann (18261866), German mathematician of genius, brought extraordinary results in almost all
fields: differential geometry (riemannian manifolds, riemannian
geometry), complex analysis (Riemann surfaces, Riemanns mapping theorem, the Riemann zeta function in relation to number
theory, the nontrivial zeros of which are famously conjectured
to lie on the line of real part 21 ), integration (Riemann integral),
Fourier analysis (Riemann-Lebesgue theorem), differential calculus (Green-Riemann theorem), series,... He also contributed to
numerous physical problems (heat and light propagation, magnetism, hydrodynamics, acoustics, etc.).
unsuccessful, by a larger domain which is simply connected.

Denote by H the helix-shaped surface which is constructed as follows.
Take an infinite collection of planes, stacked vertically. Each of these planes
is then cut, for instance along the positive real axis, and the upper part of the
cut of one plane is then glued to the the lower part of the cut in the plane
immediately above:
cutting...
... and gluing
(The third figure is what is called a helicoid in geometry. Here, only the
topology of the surface concerns us. The helicoid can be found in spiral
staircases. Close to the axis of the staircase, the steps are very steep; this is why
masons use an axis with nonzero radius.)
This surface is called the Riemann surface of the logarithm function.
Each of the points on this surface can be mapped to a point on the complex
plane, simply by projecting vertically. Two points of the helicoid precisely
vertical to each other are therefore often identified to the same complex number.
They are distinguished only by being on a different floor. Note that walking
on the helicoid and turning once around the axis (i.e., the origin), means that
the floor changes, as in a parking garage.
PROPOSITION 5.8 H is simply connected.
It it then possible to define a logarithm function on H by integrating

the function 7 1/ along a path on H joining 1 to z (1/ denoting
the inverse of the complex number identified with a point of the Riemann
surface). It is easy to see that this path integral makes sense and, because
of this proposition, it does not depend on the choice of the path from 1
Harmonic functions
139
to z. This gives a well-defined logarithm function L : H C such that

exp(L(z)) = z for all points in H .
The point of view of a physicist is somewhat different. Instead of considering H , she is happy with C \ {0}, but says that the logarithm function is
multivalued, that is, for a given z C \ {0}, the function L(z) takes not a
single value, but rather a whole sequence of values; the one that should be
chosen depends on the path taken to reach the point z. The function L is
then holomorphic on any simply connected subset of C\{0} (and the classical
local results, such as the Taylor formulas, are applicable). It is also said that
the logarithm function has many branches (corresponding to a choice of the
floor of H ).
This viewpoint can seem quite shocking to the adepts of mathematical
rigor, but it has its advantages and is particularly intuitive.
Next, let us construct the Riemann surface associated to the square-root
function. It is clear that by extending the square-root function by continuity
on C \ {0}, the sign of the function changes after one turn around the origin,
but comes back to its initial value after turning twice. Hence, we take two
copies of the complex plane, and cut them before gluing again in the following
manner:
gluing
When the square-root function has a given value on one of the sheets, it
has the opposite value on the equivalent point of the other sheet. However, it
should be noticed that this surface is not simply connected.
5.2
Harmonic functions
5.2.a
Definitions
DEFINITION 5.9 (Harmonic function) A function u : R2 R, defined on an
open set of the plane, is harmonic on if its laplacian is identically zero,

that is, if it satisfies
u =
2u 2u
+
=0
x 2 y2
at any point of .
Complex Analysis II
140
We will see, in this section, that a number of theorems proved for holomorphic functions remain true for harmonic functions, notably the mean value
theorem, and the maximum principle. This will have important consequences
in physics, where there is an abundance of harmonic functions.
First, what is the link between holomorphic functions and harmonic functions?
If f H () is a holomorphic function, it can be expressed in the form
f = u + iv, where u and v are real-valued functions; it is then easy to show
that u and v are harmonic. Indeed, f being C , the functions u and v are
also C , so
u =
u u Cauchy v v Schwarz
2u 2u
+
=
+
=
= 0
x 2 y2
x x y y
x y y x
from the Cauchy-Riemann equations on the one hand, and Schwarzs theorem
concerning mixed derivatives of functions of class C 2 on the other hand. The
argument for v is similar. To summarize:
THEOREM 5.10 The real and the imaginary part of a holomorphic function are
harmonic.
Now, we introduce the notion of harmonicity for a complex-valued function. If f is C 2 we have

2 f

2 f
f =
i
+i
f =4
f = 4 f .
+ 2 =
2
x
y
x
y
z z
x
y
DEFINITION 5.11 A complex-valued function f defined on an open subset
of the plane is harmonic if it satisfies

f = 0,
5.2.b
i.e.,
f = 0.
Properties
PROPOSITION 5.12 A complex-valued function f is harmonic if and only if its
real and imaginary parts are harmonic.
Proof. If u and v are harmonic, then u + iv is certainly also harmonic since the
operators and are C-linear.
Conversely, if f is harmonic, it suffices to take the real part of the relation f = 0
to obtain u = 0, and similarly for the imaginary part.
Hence, in particular, any holomorphic function is harmonic. What about

the converse? It cannot hold in general because there need be no relation
whatsoever between the real and imaginary parts of a complex-valued harmonic function, whereas there is a very subtle one for holomorphic functions
(as seen from the Cauchy-Riemann equations). However, there is some converse
for the real (or imaginary) part only:
Harmonic functions
141
THEOREM 5.13 Any real-valued harmonic funcion is locally the real part (resp. the
imaginary part) of a holomorphic function.
What is meant by locally here? Using only the methods of differential

equations, one can show that:
PROPOSITION 5.14 Let B be an open ball and f : B R be a harmonic
function. There exists a harmonic function v : B R such that F = u + iv is
holomorphic on B.
The astute reader will undoubtedly have already guessed what prevents
the extension of this theorem to any open subset : if has holes, the
function v may change value when following its extensions around the hole.
The same phenomenon will be seen in the section about analytic continuation,
and it has already been visible in Cauchys theorem. On the other hand, if u
is harmonic on a simply connected open subset, then u is indeed the real part
of a function holomorphic on the whole of .
COROLLARY 5.15 If u : C is harmonic, then it is of class C .
Proof. It is enough to show that for any z , there is some open ball B containing
z such that u restricted to B is C . But on B there is a holomorphic function f such
that u = Re( f ) on B; in particular, f is C on B, hence so is its real part in the sense
of R2 -differentiability.
The following theorems are also deduced from Proposition 5.14.

THEOREM 5.16 (Mean value theorem) Let be an open subset of R2 C and
let u be harmonic on . For all z0 and for all r R such that B(z0 ; r) ,
we have
Z
1
u(z + r e i ) d,
u(z0 ) =
2 0
or, in real notation (on R2 ),
Z
1
u(x0 , y0 ) =
u(x0 + r cos , y0 + r sin ) d.
2
THEOREM 5.17 (Maximum principle) Let be a domain in C, u : R
a harmonic function and a . Then either u is constant on or a is not a

local maximum of u, that is, any neighborhood of a contains a point b such that
u(b) < u(a).
Similarly, if u is not constant, it does not have a local minimum either.
For harmonic functions, it is possible to explain the maximum principle
graphically. Indeed, if z is a local extremum (minimum or maximum), and is
not constant, the function u must have the following aspect:
142
Complex Analysis II
y
y
or
But, in the first case, it is transparent that 2 u/ x 2 must be positive, as

well as 2 u/ y 2, and so the laplacian of u cannot be zero. In the second case,
those same partial derivatives are negative, and the laplacian cannot be zero
either.
A harmonic function must have one of the two partial derivatives negative
and the other positive, which gives the graph the shape of a saddle or a
mountain pass or a potato chip3:
y
x
Such a surface obviously has no local extremum, as any bicyclist will tell you!
5.2.c A trick to find f knowing u

Consider a function u : R which is harmonic on a simply connected
open set . To find a function v : R which is also harmonic and such
that u + iv is holomorphic, we have to solve the system
u
v
x = y
for any z .
v
u
= .
y
x
There is a classical method: integrating the second equation, we write
Zx
u
v(x, 0) =
(x , 0) dx ,
y
0
then
v(x, y) = v(x, 0) +
3
y
0
u
(x, y ) dy
x
A question for physicists: why is it that a potato chip has a hyperbolic geometry? Another
question: why is it also true of the leaves of certain plants, and more spectacularly, for many
varieties of iris flowers?
Harmonic functions
143
(integrating the first equation).

In fact, a trick exists to find easily the holomorphic function f with
real part equal to a given harmonic function u, avoiding the approach with
differential equations. Therefore let u be a harmonic function on a simply
connected open set C. We are looking for a holomorphic function
f : C such that u = Re( f ). Hence in particular
u(x, y) =
f (z) + f (z)
2
with z = x + i y
namely,
f (x + i y) + f (x + i y)
()
2
We are considering the function u(x, y) as a function of two real variables.
Let us extend it to a function of two complex variables e
u (z, z ). For instance,
if u(x, y) = 3x + 5x y, we define e
u (z, z ) = 3z + 5z z . Now put
z z
def
g(z) = 2e
.
()
u ,
2 2i
Then, using formulas () and (), we get
z z 2 n z
z
z
z o
g(z) = 2e
u ,
=
f
+i
+f
i
2 2i
2
2
2i
2
2i
u(x, y) =
= f (z) + f (0),
(5.3)
so that f (z) is, up to a purely imaginary constant, g(z) g(0)/2, and hence
f (z) = g(z) u(0, 0) + i Cnt .
To establish formula (5.3), we used the fact that the function z 7 f (z) can
also be expressed as a certain function fe of the variable z , and more precisely
as f (z ). For instance, if f (z) = 3z + iz 2 , then f (z) = 3z iz 2 = f (z ).
Now, isnt all this a little bit fishy? We have considered in the course of
this reasoning, that the functions (u, f ,...) behaved as operators performing
certain operations on their variables (raise it to a power, multiply by a constant, take the cosine...), and that they could just as well perform them on a
complex variable if they could do so on a real one, or on z if they knew how
to operate on z. This is not at all the usual way to see a function! However, the
expansions of f and u in power series give a rigorous justification of all those
manipulations. They also give the limit of the method: it is necessary that the
functions have power series expansions valid on the whole set considered.
Thus, this technique works if is an open ball (or a connected and simply
connected open subset, by analytic continuation), but not on a domain which
is not simply connected, an annulus, for example.4
4
I wish to thank Michael Kiessling for teaching me this little-known trick.
Complex Analysis II
144
Example 5.18 Let u(x, y) = 3x 2 3 y 2 + 5x y + 2, which satisfies u = 0. Then we have
e
u (z, z ) = 3z 2 3z 2 + 5zz + 2 and 2e
u (z/2, z/2i) = 3z 2 + 5z 2 /2i + 4 = f (z) + f (0), which gives
in particular that 2Re f (0) = 4, hence f (0) = 2+iCnt . Hence we have f (z) = 3z 2 +5z 2 /2i+2
(up to a purely imaginary constant) and one checks that the real part of f is indeed equal to
u.
Example 5.19 Consider the function u(x, y) = e x sin y. It is easy to see that this function
is harmonic. Then we have g(z) = 2e z/2 sin(z/2i) = 2ie z/2 sinh(z/2). Since g(0) = 0, we
have f (z) = g(z) + i Cnt , hence
h
z i
.
e x sin y = Re 2ie z/2 sinh
2
5.3
Analytic continuation
Assume a function f is known on an open subset by a power series
expansion that converges on . We would like to know it is possible to extend
the definition of f to a larger set.
Take a convenient example: we define the function f on the open ball
= B(0 ; 1) by
X
f (x) =
zn.
(5.4)
n=0
We know that this power series is indeed convergent on , but divergent for
|z| > 1. However, if we start from the point z0 = 12 + 2i3 , for instance, we can
find a new power series which converges to f on a certain neighborhood
of z0 :
X
f (z) =
an (z z0 )n
on a neighborhood of z0 .
(5.5)
n=0
But surprise! this neighborhood may well extend beyond , as is shown

in Figure 5.2 on the facing page.
We can then choose a path and try to find, for any point z of , a disk
of convergence of a power series centered at z coinciding with f on . This
is what is called an analytic continuation along a path (see Figure 5.3).
This means that it has been possible to continue the function f to certain
points of the complex plane where the original series (5.4) was not convergent.
On the intersection , the two series (5.4) and (5.5) coincide, because of
the rigidity of holomorphic functions.
Is this rigidity sufficient to ensure that this analytic continuation is unique?
Well, yes and no.
If we try to continue a function analytically by going around a singularity,
the continuation thus obtained may depend on which side of the singularity
Analytic continuation
z0
145
Fig. 5.2 Analytic continuation of f on an open set extending beyond .
Fig. 5.3 Continuation along a path with successive disks, and second continuation
along a path .
146
Complex Analysis II
the chosen
For instance, in the case of the function f defined by
P path takes.
n+1 z n /n, which is equal to z 7 log(1 + z) on = B (0 ; 1),
f (z) =
(1)
n=1
there exists a singularity at the point z = 1. Two continuations, defined by
paths running above or under the point z = 1, will differ by 2i.
On the other hand, the following uniqueness result holds:
THEOREM 5.20 Let f be aSfunction holomorphic on a domain 0 . Let 1 , . . . , n
n
be domains in C such that i=1 i is connected and simply

connected and conS
tains 0 . If there exists an analytic continuation of f to ni=1 i , then this analytic
continuation is unique.
5.4
Singularities at infinity
It is possible to compactify the complex plane by adding a single point at
infinity.
def
DEFINITION 5.21 (Riemann sphere) We denote by C = C {} the set
consisting of the complex plane together with an additional point infinity.
The point at infinity corresponds to the limit |z| + in C. Since

we add a single point, independent of the argument of z, this means that
all the directions to infinity are identified. Then we are led to represent
this extended complex plane in the form of a sphere,5 called the Riemann
sphere. This can also be obtained, for example, by stereographic projection
(see Figure 5.4 on the next page).
DEFINITION 5.22 (Residue at infinity) Let f be a function meromorphic on
the complex plane. Define a function F by F (z) = f (1/z). The function f
is said to have a pole at infinity if F has a pole at 0, and f is said to have an
essential singularity at infinity if F has an essential singularity at 0.
The residue at infinity of f , denoted Res ( f ; ), is the quantity
1
Res ( f ; ) = Res 2 F (z) ; 0 .
z
5
One can understand why a plane with all directions to infinity identified is a sphere by
taking a round tablecloth and grasping in one hand the whole circumference (the points of
which are infinitely far from the center for a physicist). One obtains a pouch, which is akin
to a sphere.
The open sets in C are the usual open sets in C and the complements of compact subsets
in C. The topology obtained is the usual topology for the sphere, in which the point is like
any other point.
Singularities at infinity
147
N
z
z
Fig. 5.4 The Riemann sphere, obtained by stereographic projection. If z is a point in
the complex plane C, the line (N z) meets the sphere centered at 0 with radius 1
in a single point z , the image of z on the Riemann sphere. The equator of
the sphere is the image of the circle centered at the origin with radius 1. The
point N corresponds to the point at infinity (and to going to infinity in all
directions).
The entire functions which have a pole at infinity are (only) the polynomials.
Remark 5.23 The definition of the residue at infinity is justified by the fact that it is not really
the function f which matters for the residue but rather the differential form f (z) dz (which
is the natural object occuring in the theory of path integration; see Chapter 17). By putting
w = 1/z, it follows that

1
1
1
f (z) dz = 2 f
dw = 2 F (w) dw.
w
w
w
The residue theorem can then be extended to the Riemann sphere:

THEOREM 5.24 (Generalized residue theorem) Let be an open subset of C
(the topology of which is that of a sphere in space by stereographic projection) and
K a closed subset with oriented boundary = which is made of one or more
closed paths in C (hence not passing through the point at infinity). Let f be a
function holomorphic on , except possibly at finitely many points. Denote by zk the
singularities of f inside K ; then
Z
X
f (z) dz = 2i
Res ( f ; zk ),
where the sum ranges over the singularities, including possibly the point at infinity.
Proof. It suffices to treat the case where the point at infinity is in the interior of K .
Assume that all the finite singularities (in C) are contained in the ball B(0 ; R). The
set {|z| > R} is then a neighborhood of , which we may assume is contained
in K

(increasing R if need be). We then define K = K \ z C ; |z| > R , which has a
boundary given by
K = C (0 ; R)
(with positive orientation). Applying the residue theorem to K , we find
Z
Z
X
f (z) dz +
f (z) dz = 2i
Res ( f ; zk ).
|z|=R
finite
singularities
Complex Analysis II
148
Thus we obtain
Z
f (z) dz =
|z|=R
|w|=1/R

1
1
f
dw = 2i Res ( f ; ).
w2
w
(Note that the change of variable w = 1/z transforms a positively oriented contour to
a negatively oriented one.)
The stated formula follows.
Example 5.25 The function f : z 7 1/z has a pole at infinity, with residue equal to 1.
If we let = C (5 ; 1), a contour which does not contain any singularity in C, we have
Z
f (z) dz = 0.
Using the generalized residue theorem , we get

Z
h
i
f (z) dz = 2i Res ( f ; 0) + Res ( f ; ) = 2i(1 1) = 0.
5.5
The saddle point method
The saddle point method is a computational recipe developed by Debye7
to compute integrals of the type
Z
def
I () = e f (z) dz.
Here, z is a real or complex variable.8 The goal is to find an approximation

of I () in the limit when goes to +, assuming that the function f admits
a local minimum.
6
With respect to the closed set K which is the outside of on the picture. In fact, the
notions of the parts inside and outside a curve do not make sense on the Riemann sphere.
Which side is the outside of the equator on Earth? The northern hemisphere or the southern
hemisphere?
7
Petrus Josephus Wilhelmus Debye (18841966), a Dutch physicist, born in Maastricht, studied the specific heat of solids (Debye model), ionic solutions (Debye potential, see problem 5
p. 325, Debye-Hckel method), among other things. He received the Nobel prize in chemistry
in 1936.
8
There exist generalizations of this method for functions of vector variables, even for fields,
which are used in quantum field theory.
5.5.a
149
The general saddle point method
Consider, more generally, an integral of the type

Z
def
I = e f (z) dz,
(5.6)
where is a complex parameter and is a path in C. We assume that f is

a function analytic on an open subset containing , and we want to evaluate
this integral in the limit where || .
The first idea is to evaluate the integral in the parts of the domain of
integration where the real part of f is minimal, since it is in those regions
that the modulus of exp f is largest and so they should bring the main
contribution to the integral. Unfortunately, in those regions, the imaginary
part of f may very well vary and, therefore, e f will oscillate, possibly very
rapidly. These oscillations may well destroy a large part of the integral by
compensations, and the suggested method will not be justified.
The idea of Debye is then to deform the contour into an invented path
satisfying the following properties:

i) along , Im f (z) is constant;
ii) there exists a point z on such that

d f (z)
= 0;
dz z
iii) the real part of f has a local minimum at z = z ;

iv) moreover, this minimum becomes steeper and steeper;
v) finally, the initial path can be deformed continuously into (i.e.,
during the deformation, it must not pass through any of the singularities of f ).
DEFINITION 5.26 A point of C satisfying Conditions ii) and iii) above is
called a saddle point for f .

The condition Im f (z) = Cnt gives, in general, the equation for the
curve . Under this condition, the phase of the exponential is constant, and
the computation reduces to a real integral.
Consider now the behavior of the function f around the saddle point, if
it exists. Suppose therefore that

d f
d2 f
=0
and
6= 0.
dz z=z
dz 2 z=z
Complex Analysis II
150
Then we can write
Denoting

d2 f
1
2
f (z) = f (z ) + (z z )
+ o(z z )2 .
2
dz 2 z=z

d2 f
= e i
and
z z = r e i ,
dz 2 z=z
we have, as a first approximation,

Re f (z) f (z ) = 12 r 2 cos(2 + ) + o(r 2 ),

Im f (z) f (z ) = 12 r 2 sin(2 + ) + o(r 2 ).
(5.7)
Hence we see that, at the saddle point z , the angle between the path
and the real axis should be such that sin(2 + ) = 0, which leaves two
possible directions, orthogonal to each other (solid lines in Figure 5.5). The
twodirections characterized
by cos(2 + ) = 0 are those where the real part

of f (z) f (z ) is constant. Between those lines, the real part is either
positive or negative, depending on the sign of cos(2 + ), which shows that
we are indeed dealing with a saddle point, or a mountain pass (see Figure 5.5).
The two orthogonal directions where the imaginary part of f (z) f (z )
vanishes (to second order) are precisely the directions where the real part of
f (z) f (z ) varies most quickly, either increasing or decreasing.
Consequently, we must choose, at the saddle point, the direction characterized by sin(2 + ) = 0 and such that, along this direction, Re( f ) has a
local minimum at z , and not a local maximum.
We have seen that, because of the properties of holomorphic functions, it
is usually possible to find a path along which the oscillations of e f have disappeared and where, moreover, f will have a local minimum as pronounced
as possible. This justifies the other common name of the method: the method
of steepest descent.
There only remains to evaluate the integral along the path . We introduce a new variable t parameterizing the tangent line to at z , so that we
can write
t2
f (z) = f (z ) + + o(z z )2 ,
2
and, from (5.7), we have
d2 f
t = (z z )
dz 2
(the sign of the square root being chosen in such a way that its real part is
positive). We now rewrite formula (5.6) in terms of this parameter t, which
gives, noticing that dz = dz(t)/dt dt,
Z +
2 dz(t)
I e f (z )
e t
dt.
dt
151
Re( f )
d
a
c
b
Fig. 5.5 The two solid curves are perpendicular and locally determined by the condition
sin(2 + ) = 0. These are not only the curves such that Im( f ) is constant,
but also those where Re( f ) varies most rapidly. The two dotted lines separate
the regions a, b , c and d . In regions a and c , Re f (z) > Re f (z ), whereas in
regions b and d , Re f (z) < Re f (z ). The chosen path therefore goes through
regions a and c .
Since the function d f /dz varies slowly in general, compared to the gaussian
term, we can consider that dz/dt = (d2 f /dz 2 )1/2 is constant, and evaluate
the remaining gaussian integral, which gives
I e f (z )
2

2
d f /dz 2 z=z
)1/2
(5.8)
This formula can be rigorously justified, case by case, by computing the nongaussian contributions to the integral and comparing them with the main
term we just computed. The computation must be done by expanding f to
higher order. One is then led to compute gaussian moments (which is easy) to
derive an asymptotic expansion of I to higher order. A condition to ensure
that the higher orders do not interfere is that the ratio between
p
1/ p
2
1/2
d f
d f
and
dx p x=x
dx p x=x
tends to 0 when tends to infinity. This condition is necessary, but not
sufficient.
Complex Analysis II
152
5.5.b
The real saddle point method
There exists a particularly simple case of the problem when f is a real-valued

function. Then it is not possible to work in the complex plane, but a gaussian
integral computation leads to a similar result.
Consider an integral of the shape
Z
I =
e f (x) dx,
R
where f is a function at least twice continuously differentiable, depending on

a parameter , and with the following property:
for every , f has a unique minimum9 at a point x , de

noted m = f (x ); moreover, this minimum is more and

(5.9)
def
2
2
more narrow; i.e., if we denote M = d f /dx x=x , we

have M > 0 and M + as .
The dominating contribution to the integral is evidently given by the

region of x close to the local minimum. So the idea is to expand the function
f around this minimum. If it suffices to expand to order 2 (the first derivative
is zero because the point is a local minimum), then the computation reduces
to a gaussian integral. This approximation should be all the more justified as
we take the limit [ ].
We thus write
1
f (x) = m + M (x x )2 + o(x x )2 ,
2
and as a first approximation estimate the integral by

Z
1
2
m
e 2 M (xx ) dx.
I e
We recognize here a gaussian integral, so the result is
I
e m
2
= e m
M
d2 f
2

.
2
/dx x=x
The first factor is simply the value of the exponential function at its maximum, and the second factor
p expresses the characteristic width of this maximum, proportional to 1/ M . This result should be compared to the result
of (5.8).
9
If f has several local minima, each one contributes to the integral.
Exercises
153
Remark 5.27 Sometimes, an integral of the same type as before arises, but with an imaginary
exponential, that is, of the type
Z
I = e i f (x) dx,
R
where f is a real-valued function.

In this case, if f (x0 ) varies fast enough as tends to infinity for a given x0 , we expect that
the neighborhood of this point will only have a small contribution to the integral, because
of destructive interferences. This phenomenon can be understood by hand by the following
rough reasoning: if we expand around x0

df
+ o(x x0 ),
f (x) = f (x0 ) + (x x0 )
dx x=x0
and if d f /dx 6= 0, we have, as a first approximation, an integral of the type

Z
def d f
M e i k(xx0 ) dx,
with k =
,
dx x=x0
which is zero if k 6= 0.
We must therefore look for extremums of f to find important contributions, which are
called for obvious reasons stationary phase points; assuming there is a unique such point x , a
computation very similar to the previous one for the real saddle point method gives the result
1/2
2i
.
I e i f (x )
d2 f /dx 2 | x=x
This case is also called the stationary phase method.
Remark 5.28 The saddle point method is not just a trick for physicists. It is very often used
in pure mathematics, analytic number theory in particular; for instance, it occurs in the proof
of the irrationality of one of the numbers (5), (7), . . . , (21) by Tanguy Rivoal [73]! (Recall
that
X 1
(s) =
ns
n1
for Re(s) > 1).
EXERCISES
Exercise 5.1 Define
u(x, y) = e x (x cos y + y sin y).

Find so that u is the real part of a holomorphic function f (identifying the plane R2 and
the complex plane). What is the function f ?
Exercise 5.2 Similarly, find all the holomorphic functions with real part equal to
P (x, y) =
x(1 + x) + y 2
.
(1 + x)2 + y 2
Exercise 5.3 Show that, in two dimensions, a cavity in a conductor, devoid of charge, has
potential identically zero.
154
Complex Analysis II
Exercise 5.4 (Stirling formula) Using the Euler function defined on the complex halfplane Re(z) > 0 by the integral representation
Z +
def
(z) =
e x x z1 dx,
0
and using the fact that (n + 1) = n! for all n N, show, by the complex saddle point method,
that
p
(n + 1) = n! nn e n 2n.
This formula, established by Abraham de Moivre, is called the Stirling formula.
SOLUTIONS
Solution of exercise 5.3. Denote by the domain defined by the cavity; is simply con-
nected. The electrosatic potential satisfies (z) = 0 at any point z . This potential is also
a harmonic function, hence achieves its minimum and its maximum on the boundary .
Consequently, we have (z) = 0 for any z .
Solution of exercise 5.4. The first thing to do is to find a suitable contour. To start with,
write
( + 1) =
e x log x dx,
then look for the points where d f /dz = 0, with f (z) = z log z. We find z = , and
d2 f /dz 2 |z = 1/, which is a positive real number. The two lines for which, locally, the
imaginary part of f is constant are the real axis and the line with equation z = + i y. It is
the first which corresponds to a local minimum of f at z = ; therefore the real axis is the
most suitable contour.
Hence we evaluate the integral according to the formula (5.8), which gives
p
( + 1) e log 2
[ ].
One can check explicitly that the non-gaussian contributions (those given by the terms of order
larger than 2 in the expansion of f ) are inverse powers of . Hence
p

( + 1) = e 2 1 + O (1/) ,
and in particular, for integral values of ,
p
(n + 1) = n! nn e n 2n.
Notice that this approximation is very quickly quite precise; for n = 5, the exact value is
5! = 120 and the Stirling formula gives approximately 118, which is an error of about 2%.
Even more impressive, the relative error from applying the Stirling formula for n = 100 is
about 0, 08%. . . (note that this error is nevertheless very large in absolute value: it is roughly of
size 10155 , whereas 100! is about 10158 ).
Remark 5.29 If one seeks the asymptotic behavior of (z) for large z not necessarily real, the
computation is more or less the same, except that the path is not the real axis any more.
However, the result of the computation is identical (the result is valid uniformly in any region
where the argument ], ] of z satisfies + < < for some > 0).
Chapter
Conformal maps
6.1
Conformal maps
6.1.a
Preliminaries

Consider a change of variable f : (x, y) 7 (u, v) = u(x, y), v(x, y) in the
plane R2 , identified with C. This change of variable really only deserves the
name if f is locally bijective (i.e., one-to-one); this is the case if the jacobian
of the map is nonzero (then so is the jacobian of the inverse map):

u u
x x

D (u, v) x y
D (x, y) u v

=
6= 0

and
D (u, v) = y y 6= 0.
D (x, y) v v

x y
u v
THEOREM 6.1 In a complex change of variable
z = x + i y 7 w = f (z) = u + iv,
and if f is holomorphic, then the jacobian of the map is equal to

D (u, v) 2
= f (z) .
J f (z) =
D (x, y)
Conformal maps
156
Proof. Indeed, we have f (z) =

relations,
2
f (z) =
u
x
2
v
u
+i
and hence, by the Cauchy-Riemann
x
x
v
x
2
u v
v u
= J f (z).
x y x y
DEFINITION 6.2 A conformal map or conformal transformation of an

open subset R2 into another open subset R2 is any map f :
, locally bijective, that preserves angles and orientation.
THEOREM 6.3 Any conformal map is given by a holomorphic function f such that
the derivative of f does not vanish.
This justifies the next definition:

DEFINITION 6.4 A conformal transformation or conformal map of an
open subset C into another open subset C is any holomorphic

function f : such that f (z) 6= 0 for all z .
Proof that the definitions are equivalent. We will denote in general w = f (z).
Consider, in the complex plane, two line segments 1 and 2 contained inside the set
where f is defined, and intersecting at a point z0 in . Denote by 1 and 2 their
images by f .
We want to show that if the angle between 1 and 2 is equal to , then the same
holds for their images, which means that the angle between the tangent lines to 1 and
2 at w0 = f (z0 ) is also equal to .
Consider a point z 1 close to z0 . Its image w = f (z) satisfies
lim
zz0
and hence
w w0
= f (z0 ),
z z0
lim Arg(w w0 ) Arg(z z0 ) = Arg f (z0 ),
zz0
which shows that the angle between the curve 1 and the real axis is equal to the angle
between the original segment 1 and the real axis, plus the angle = Arg f (z0 ) (which
is well defined because f (z) 6= 0).
Similarly, the angle between the image curve 2 and the real axis is equal to that
between the segment 2 and the real axis, plus the same .
Therefore, the angle between the two image curves is the same as that between the
two line segments, namely, .
Another way to see this is as follows: the tangent vectors of the curves are transformed according to the rule v = d f z0 v. But the differential of f (when f is seen as
a map from R2 to R2 ) is of the form
P
P
x
y cos sin

,
(6.1)
d f z0 =
Q Q = f (z0 ) sin
cos
x
y
where is the argument of f (z0 ). This is the matrix of a rotation composed with a
homothety, that is, a similitude.
Conformal maps
157
1
1
2
Conversely, if f is a map which is R2 -differentiable and preserves angles, then at

any point d f is an endomorphism of R2 which preserves angles. Since f also preserves
orientation, its determinant is positive, so d f is a similitude, and its matrix is exactly
as in equation (6.1). The Cauchy-Riemann equations are immediate consequences.
Remark 6.5 An antiholomorphic map also preserves angles, but it reverses the orientation.
6.1.b
The Riemann mapping theorem
DEFINITION 6.6 Two open subsets U and V of C are homeomorphic if
there exists a bijection f : U V such that f and its inverse f 1 are both
continuous. Such a map is called a homeomorphism of U into V .
Remark 6.7 In fact, in this case the continuity of f 1 is an unnecessary condition, because it
turns out to be implied by the continuity of f , the fact that f is bijective, and that U and V
are open. However, this is quite a delicate result.
Example 6.8 The complex plane C is homeomorphic to the (interior of the) unit disc, through
the homeomorphism
f : z 7
See also Exercise 6.2 on page 174.
z
.
1 + |z|
DEFINITION 6.9 Two open sets U and V are conformally equivalent if there
exists an injective conformal application f : U C such that V = f (U ).
THEOREM 6.10 (Riemann mapping theorem) Let U be any open subset of C

such that U 6= C. If U is homeomorphic to the unit disc D = B(0 ; 1), then U is
conformally equivalent to D .
In general, two open sets which are homeomorphic are not conformally
equivalent. For instance, two annuli r1 < |z| < r2 and r1 < |z| < r2 are
conformally equivalent if and only if r2 /r1 = r2 /r1 .
So the Riemann mapping theorem is a very strong result, since any connected and simply connected open set, except C, can be transformed, via a conformal map, into the unit disc.
To illustrate, consider in Figure 6.1 on the next page some open subset
U of the complex plane. Riemanns theorem guarantees the existence of a
conformal map w = f (z) such that B(0 ; 1) = f (U ). The boundary U
of U is mapped to the unit disc C (0 ; 1). (This seems natural, but to make
Conformal maps
158
z
U
w = f (z)
D
Q
0
Fig. 6.1 The open subset U being homeomorphic to the unit disc, it is conformally
equivalent to it. There exists therefore a conformal map w = f (z) such that the
image of U by f is D. One can ask that f (P ) = 0 and f ( Q ) = Q .
sense of it we need to know in general that f extends to the boundary of U .

That this is the case is shown by Carathodory and Osgood, provided this
boundary is a simple closed continuous curve; in the explicit examples below,
this extension will be obvious.) Moreover, it is possible to impose additional
conditions on f , for instance, by asking that a particular point P U be
mapped to the center of the unit disc: f (P ) = 0, and by asking also that a
point Q U be mapped to a given point on the unit circle.
It is important to note that, while Riemanns theorem tells us that the
conformal representation of U exists theoretically, it does not tell how to find
it. Physicists like Joukovski1 found clever conformal representations of various
domains in the plane; these are useful to compute easily (in terms of computer
time) the effects of a fluid flowing around those domains.
6.1.c Examples of conformal maps

The following examples of conformal maps are given without proof. The
reader is invited to check their properties by herself.
Representation of a half-plane into a disc
Let z0 be any point in the open upper half-plane. Let
w = f (z) =
z z0
.
z z0
The upper half-plane is then transformed bijectively into the interior B(0 ; 1)
of the unit circle. The point z0 is mapped to 0. The real points x tending
to either + or are sent to the same limit 1. As z runs over the real
1
Nicola Joukovski (18471921), a Russian physicist, studied hydrodynamics and aerodynamics. He discovered a conformal map z 7 21 (z + a 2 /z), which transforms the circles going
through a and a into a profile similar to an airplane wing.
Conformal maps
159
axis in increasing order, w goes around the unit circle in the counterclockwise
direction, starting and ending at the point w = 1.
z
w = f (z)
P (z0 )
B
P (0)
C
C
Infinite strip of width a

Using the complex exponential function, we can transform an infinite strip
of width a into the upper half-plane, by putting
w = f (z) = e z/a .
The points on the real axis and the line Im(z) = a in the limit Re(z)
both converge to 0, while in the limit Re(z) +, the real points converge
to + and those with height a converge to (on the real axis).
D
F
z
w = f (z)
1
F E D
1
A B
Half-plane outside a half-circle

The upper half-plane minus a half-circle is mapped to the upper half-plane by
putting
a
1
w = f (z) =
z+
,
2
z
where a is an arbitrary real number. This transformation is not conformal at
the points 1 and 1 (the derivative vanishes), which explains that the right
angles between the half-circle and the real axis are transformed into flat
angles (equal to ).
z
w = f (z)
A B
0
C
Conformal maps
160
Half-circle
The unit upper half-circle is mapped into the upper half-plane by putting
1+z 2
w = f (z) =
.
1z
The point 1 (C in the figure) is mapped to 0 while the point 1 (A) is mapped
to infinity on the real axis (either + or , depending on whether one
approaches A on the real axis or on the boundary of the half-circle).
z
w = f (z)
1
1
C
Some conformal mappings into the unit disc

The exterior of the unit disc is mapped conformally to the interior by simply
putting w = 1/z.
Also the exterior of an ellipse with major axis equal to 2 cosh() and minor
axis equal to 2 sinh() is mapped to the interior of the unit disc by putting
1
1
w = f (z) =
z e + e .
2
z
Other examples
The angular sector {r e i ; 0 < r < 1, 0 < < 2 } (a quarter-disc) is mapped
into the upper half-disc by putting w = z 2 .
Similarly, the angular sector of angle /k (with k R, k 12 ) is mapped
into the upper half-plane by putting w = z k .
The upper half-disc in mapped to a semi-infinite strip (infinite to the left)
by putting w = log z with the logarithm defined by a cut in the lower halfplane.
D
z
w = f (z)
1
B
0
A
The same transformation maps the upper half-plane into an infinite strip.
Conformal maps
161
The German mathematician Hermann Schwarz (18431921) succeeded his teacher Weierstrass at the University of Berlin in 1892.
He worked on conformal mappings, potential theory (proving
the existence of a solution to the Dirichlet problem in 1870), and
partial differential equations.
6.1.d The Schwarz-Christoffel transformation

We are now looking for a conformal transformation that will map a domain
bounded by a polygonal line into a simpler domain, for instance, the upper
half-plane (or the unit disc, but we already know how to go from the unit disc
to the upper half-plane). For this, we use the following result:
THEOREM 6.11 (Schwarz-Christoffel transformation) Let n N be an integer,
x1 , . . . , xn real numbers, and 1 , . . . , n real numbers between 0 and 2. Let finally

a C be an arbitrary complex number. The map given by
dw
= a(z x1 )1 /1 (z x2 )2 /1 (z xn )n /1
(6.2)
dz
transforms the real axis into a polygonal curve with successive angles (1 , . . . , n ) and
transforms the upper half-plane into the domain bounded by this curve (see Figure 6.2
on the following page).
Formula (6.2) provides only a local expression of the map. By integrating
along a complex path, we obtain the equivalent formula
Z Y
n
i
w = f (z) = a
(z x i ) 1 dz + w0 ,
(z) i=1
where w0 C is an arbitrary constant and (z) is a path coming from infinity,

contained in the open upper half-plane and ending at z. Since the upper halfplane is simply connected and the only poles of the function to integrate are
on the real axis, the result does not depend on the actual path we choose.
Remark 6.12
Changing w0 amounts simply to translating the image of the map.

Changing the value of the constant a gives the possibility to make rotations (by changing the argument of a) or homotheties (by changing its module).
By imposing the condition

X i
1 = 2,
Conformal maps
162
w = f (z)
A
A
C D
x1 x2
x3 x4
x5
D
2
B
Fig. 6.2 The Schwarz-Christoffel transform maps the open upper half-plane into the
domain bounded by the polygonal curve with successive angles 1 , . . . , n .
the polygonal line is closed and we obtain a closed polygon. An infinite polygon (when
the condition does not hold) is a limiting case of closed polygons.
If the problem to solve is to find a conformal mapping transforming the upper halfplane into a polygonal domain (either closed or open), it is possible to choose arbitrarily
three of the real numbers (x1 , . . . , xn ). The others are then imposed by the various
quantities in the picture.
It is often convenient to move one of the xi (for instance, x1 or xn ) to infinity, which
simplifies the computations.
Example 6.13 We describe how to use the Schwarz-Christoffel transform to map conformally
an infinite half-strip into the upper half-plane.
w = f (z)
1
b
The two angles in the polygonal line are equal to

x2 = 1. So we need to solve
and we fix the choice of x1 = 1 and
a
dw
1
1
= a(z + 1) 2 (z 1) 2 = p 2
=p
,
dz
z 1
1 z2
where a is a constant to be determined and = ia. We therefore put
w = f (z) = arcsin(z) + B,
and to have f (1) = b and f (1) = b , we must take B = 0 and = 2b/. Hence the
solution is
w
2b
w=
arcsin(z),
i.e.,
z = sin
.
2b
Applications to potential theory
163
Elwin Christoffel (18291900), born in Montjoie Aachen (present-day Monshau) in Germany, from a family of clothes merchants, studied in Berlin from 1850, where he had in particular
Dirichlet as a teacher (see page 171). He defended his doctoral on
the flow of electricity in homogeneous bodies, then taught there
before ending his career in the new University of Strasbourg.
Shy, irritable, asocial, Christoffel was nevertheless a brilliant
mathematician who studied Riemannian geometry, introducing
the Christoffel symbols, as well as conformal mappings.
6.2
(This section may be omitted for a first reading; it is preferable to study it with some
knowledge of Green functions, which we will present later.)
Conformal transformations are useful to solve a number of problems of
the type Poisson equation with boundary condition, for instance, the problem of the diffusion of heat in a solid, in a stationary regime, with the
temperature fixed on the boundary of the solid; or the problem of the electrostatic field created by a charge distribution in the presence of conductors.
The main idea is to transform the boundary of the problem, which may
be complicated, into a simpler boundary, by a conformal map. In particular,
the Dirichlet problem for a half-plane or for a circle (see page 170) can be
solved, and this makes it possible to solve it, at least in principle, for more
complicated boundaries.
THEOREM 6.14 Let R2 and R2 be two domains of R2 C and f a
bijective conformal map from into . Let : R be a real-valued function,
2
2
def
with laplacian denoted by (x, y) = x2 + y2 . Moreover, denote (u, v) = f (x, y)
and consider the image of by the change of variable f :

def
(u, v) = f 1 (u, v)
def
(in other words, = f 1 ). Then, if we denote (u,v) =

laplacian of in the new coordinates, we have

2
(x, y) = f (z) (u,v) .
2
u2
2
v2
the
Proof. This is an elementary but tiresome computation. It suffices to use the chain
Conformal maps
164
rule formulas
x
y
=
+
u
u x
u y
and
x
y
=
+
v
v x
v y
for partial derivatives, together with the Cauchy-Riemann relations.
COROLLARY 6.15 The property that a function is harmonic is invariant by confor-
mal transformation.
Hence, suppose we have to solve a problem like the Dirichlet problem,

namely, we must find a harmonic function on a domain domain R2
such that (z)
= f (z), where f : C is a given function. If we know
how to transform by conformal mapping into a half-circle or a half-plane,

we can solve the associated Dirichlet problem and then perform the inverse
change of variable, which will preserve the harmonicity of the solution.
However, in general the goal is to solve not a Laplace equation = 0,
but rather a Poisson equation
(x, y) = (x x0 , y y0 ).
It is therefore necessary to first determine how this equation is transformed
by a conformal mapping. In a dilation (change of scale) x 7 a x in the
argument of a Dirac distribution (see Chapter 7), a factor 1/ |a|d appears,
where d is the dimension. Here, we deal with two dimensions and, locally,
a conformal transformation is simply a similitude, with homothety factor
given by | f (z)|. We have therefore the following result (see Theorem 7.45 on
page 199 for a generalization when f is not bijective):
THEOREM 6.16 Let R2 and R2 be two domains of R2 C and let f
be a bijective conformal mapping of into . Then for (x0 , y0 ) , we have

1
f (x, y) f (x0 , y0 ) =
(x x0 , y y0 ).
f (x0 , y0 )2
As a consequence, we have the following result:
THEOREM 6.17 Let R2 and R2 be two domains of R2 C and let f
be a bijective conformal mapping of into . Let : R be a real-valued

function and let be the image of by the change of variable f . Then, if satisfies
the Poisson equation
(x, y) = (x x0 , y y0 ),
its image by the conformal map f also satisfies the Poisson equation in the new
coordinates:

(u,v) = f (x, y) f (x0 , y0 ) .
6.2.a
165
Application to electrostatics
We will give an example of application of those theorems to solve the Dirichlet

problem in the situation corresponding to electrostatics. We consider, in a twodimensional space, a conductor which leaves a hole shaped like an infinite
half-strip, denoted (see Figure 6.3).
Now, we would like to compute the potential created by an arbitrary distribution of charge (x, y). For this, because of the linearity of the electrostatic
Poisson equation ( = /0 ), we will see (Chapter 7) that it suffices to find
the potential G (x, y ; x0 , y0 ) created at the point (x, y) by a Dirac distribution
at (x0 , y0 ), which satisfies
2
2
(
G (x, y ; x0 , y0 ) = (x x0 , y y0 ) where = 2 + 2 ,
x
y
(C )
G (x, y ; x0 , y0 ) = 0 for any (x, y) and (x0 , y0 ) .
Once this function is found, the solution to the problem is given by2
ZZ
1
(x, y) =
(x , y ) G (x, y ; x , y ) dx dy .
0
One should note that the function G (x, y ; x , y ) is not invariant under translations and therefore it does not depend only on x x and y y , because
of the presence of the conductor. This function is called the Green function
for the problem.
To find this function G (x, y ; x , y ), we start by introducing complex variables z = x + i y and w = u + iv; then we perform the conformal transformation w = sin(z/2b) which, as we have seen, transforms into the upper
half-plane (see Figure 6.3). Then we look for the Green function of the new
problem, G (w ; w ), which satisfies, according to Theorem 6.17,
(u,v) G (u, v ; u , v ) = (u u , v v ),
(C )
G (u, v ; u , v ) = 0
(u, v) f ( ), (u , v ) f (),
where the image of the boundary is the real axis f ( ) = R. In other
words, we are looking for the electrostatic response of a simpler system (a
half-plane) in the presence of an electric charge.
To solve this last problem, one can use the method of virtual images (which
is justified, for instance, by means of the Fourier or Laplace transform): the
potential created by a charge at w in the neighborhood of a conductor in
the lower half-plane is the sum of the free potentials3 created, first, by the
particle at w and, second, by another particle, with the opposite charge,
placed symmetrically of the first with respect to the real axis (namely, at w ).
2
The physicists proof being that the quantity 1
(x , y ) G (x, y ; x , y ) dx dy is the
0
potential created at (x, y) by the elementary charge (x , y ) dx dy .
3
I.e., of the Green functions tending to 0 at infinity.
Conformal maps
166
z
w = f (z)
b
(a)
1
(b)
Fig. 6.3 (a) The domain left free by the conductor is the set of points z = x + i y
with y > 0 and b < x < b . The Dirichlet conditions are that the potential
is constant, for instance, zero, on . (b) The image f () of the domain is
the upper half-plane.
However, we will see in the chapter concerning distributions that the Coulomb
potential4 in two dimensions is simply

q
(w) =
log w w .
20
G is thus equal to
G (w ;
w )

w w
q
.

=
log
2
w w
The Green function of the initial problem is obtained by putting, in the last
formula, w = sin(z/2b), which gives

z0
z

sin
sin
q

2b
2b ,
G (z ; z0 ) =
log
z
z

2
sin
sin 0
2b
2b
or, by expanding in real form sin(a + ib) = sin a cosh b + i cos a sinh b,
G (x, y ; x0 , y0 ) =
q
4
x
y
x
y0 2
x
y
x
y0 2

sin
cosh
sin 0 cosh
+ cos
sinh
cos 0 sinh

2b
2b
2b
2b
2b
2b
2b
2b
log x

y
x0
y0 2
x
y
x0
y0 2 .

+
cos
sin
cosh
sin
cosh
sinh
+
cos
sinh

2b
2b
2b
2b
2b
2b
2b
2b
Which, as the reader will notice, is not properly speaking something very
intuitive. However, it is remarkable how this computation has been done so
easily and without difficulty. This is the strength of conformal mappings.
4
The Coulomb potential in any dimension is a potential satisfying the Laplace equation
= /0 ; Theorem 7.54 on page 208 shows that the three-dimensional Coulomb potential
is given by r 7 1/40 r , whereas in two dimensions Theorem 7.55 on page 208 shows that it
is r 7 (log r )/20 .
167
6.2.b Application to hydrodynamics

Consider the bidimensional flow of a perfect fluid (i.e., an incompressible, nonviscous, and irrotational fluid). Then (see, for instance, the excellent book by
Guyon et al. [44]) the fluid velocity is obtained locally, at any point, from a
potential , in the sense that the x and y coordinates of the velocity are given
by
vx =
and
vy =
x
y
This potential is harmonic since div v = 0 is equivalent to = 0. Solving
the flow problem amounts to finding with the right boundary conditions
(normal velocity equal to zero on all impenetrable surfaces). This is called the
Neumann problem.5
Remark 6.18 The velocity is obtained from a global potential if the domain under
consideration C is simply connected. Otherwise, it may not be possible to find a

global potential; it is then necessary to cut into simply connected pieces, and try,
after solving the problem on each component, to glue the solutions back together. In
what follows, we assume that is indeed simply connected.
DEFINITION 6.19 The complex velocity of the fluid is the complex quantity
given by
def
V = v x + iv y .
Since is harmonic, we know (Theorem 5.13) that there exists a harmonic
def
function : R which is conjugate to , that is, such that = + i is
holomorphic on . Then the complex derivative of is the conjugate of the
complex velocity:
(z) =
=
=
+i
z
x
x
x
Cauchy-Riemann
i
= v x iv y = V (z).
x
y
The complex velocity is therefore an antiholomorphic function.

DEFINITION 6.20 The holomorphic function is called the complex potential. A stagnation point for the flow is any point z0 where the complex
velocity vanishes: (z0 ) = 0.
Consider, for instance, the flow of a perfect fluid in a domain D in R2

(identified with the complex plane C), delimited by the ground and by a wall
5
Carl Gottfried Neumann (18321925), German mathematician and physicist, was the son
of the physicist Franz Ernst Neumann, who introduced the notion of potential in electromagnetism.
Conformal maps
168
w = f (z)
Fig. 6.4 On the left: the wind over a wall. The white region is D. On the right: the
wind over the ground, without a wall. The white region is D .
of height a (see Figure 6.4, left). We will, by a conformal transformation,

map D into the upper half-plane

def
D = z C ; Imz > 0
p
(see Figure 6.4, right). For this, it is enough to put f (z) = z 2 + a 2 = w.
Indeed, if we decompose the transformation f in three steps
p
f = (z 7 z) (z 7 z + a 2 ) (z 7 z 2 ),
p
where the function z 7 z is taken with a cut along the positive real axis,
then the domain D is transformed successively as follows:
ia
a 2
z 7 z 2
z 7 z + a 2
z 7
a 2
We can now deduce from this the flow in D with the boundary conditions6:
lim v(z) = (v , 0).
Re(z)
In the general case, one can proceed as follows. Search for the flow in the form
(z) = v z + g(z),
where g is a holomorphic function which describes a perturbation from the homogeneous flow
0 (z) = v z, such that lim|z| g (z) = 0 so that the limit conditions are respected.
169
After the conformal mapping w = f (z), the flow still obeys the same
boundary conditions, since f (z) 1 as |z| tends to infinity in the upper
half-plane. The desired flow in the upper half-plane without a wall is easy to
find: it a a uniform flow (the wind over the plain...), which can be described
by the complex potential
(w) = v w.
Coming back to the variable w:
p
p

(z) = (w) =
z 2 + a 2 = v z 2 + a 2 .
Two special points appear: the velocity of the flow is infinite on the top of
the wall, at z = ia (this is why the wind whistles when going over a pointed
obstacle: locally supersonic velocities are reached); it is zero at the foot of
the wall (around z = 0). These two points are of course the two points
where the transformation is no longer conformal7; it is therefore possible that
singularities appear in this manner, which did not appear in the flow on the
upper half-plane.
6.2.c
Potential theory, lightning rods, and percolation
From the previous example we can understand the principle of the lightning
rod. If we denote by ( r) the electrostatic potential, then close to a pointed
object (like the wall above, or an acute angle in the definition of the domain),
we can expect that will diverge. Identifying C and R2 , it is the gradient
of that will diverge, namely, the electric field. But we know that in the air
an electric field of intensity larger than some critical value, the breakdown
threshold, will cause a spontaneous electric discharge.8 Saint Elmos fire, well
known to mariners, has the same origin.
But this is not yet the most remarkable thing. Consider a problem of
random walk of a particle in an open domain C, assuming that, when the
particle meets the boundary of , it will stop. (It is possible to make numerical
simulations of this problem.) What will happen after many particles have
performed this process? One may think that the particles will accumulate on
the boundary , which will therefore grow, more or less regularly as time
goes. In fact, this is not what happens. If the boundary has cusps,
they will attract more particles and, hence, will grow faster than the rest of
the boundary. If, at a given time and due to the effect of chance, a cusp is
created (by accumulation of some particles) on a smooth part of , then
this cusp will, in turn, attract more particles. In a word: the fluctuations of
the boundary become unstable.
7
Indeed, at the foot of the wall, we have f (z) = 0, whereas around the top, f (z) diverges.
Of course, the rough explanation given only concerns lightning rods in a two-dimensional
world. In three dimensions, another argument must be found, because the one given uses
specific properties of C and not of R2 .
8
170
Conformal maps
How to explain this phenomenon? If we denote by p(z, t) the probability

density for the particle to be at position z at time t, one can show that this
probability satisfies an equation of diffusion, which is the manifestation on a
macroscopic scale of the extremely erratic behavior of microscopic particles
(the so-called Brownian motion). The diffusion equation takes the form

p(z, t) = 0
t
at any point of . We will therefore have an equation of the same type if we
perform a conformal transformation (independent of time) on the domain .
Once more, the cusps will attract more particles and we will observe a growth
phenomenon.9 (Here it is the flow of particles that diverges.)
6.3
Dirichlet problem and Poisson kernel
The Dirichlet problem is the following: the values of a harmonic function
on the boundary of a domain in the plane are assumed to be known,
and we want to find the values of the function at any point in .
For instance, one may be interested in the case of a metallic disc,10 which is
a heat conductor, in a stationary state, for which we impose a temperature T0
on one half-circle, and a temperature T1 on the other half-circle. The question
is now to determine the temperature at any point inside the disc.
More precisely, if is a domain (hence, a connected open set) of the
complex plane C, if denotes its boundary and if u0 : C is a
continuous function, we are looking for a u : C of C 2 class, such that
u
= u0 ,
u = 0 in .
A first important case is when the domain is an open ball. Indeed,

according to Riemanns mapping theorem, one can in principle reduce to this
case whenever is homeomorphic to an open ball (and not the whole of C).
In this situation, we have a first useful theorem:
THEOREM 6.21 Let u be a function harmonic on a neighborhood of B (0 ; 1). Then,
for any point a B(0 ; 1), we have

Z
1 2
1 |a|2
i
u(a) =
u(e )
d.
a e i 2
2 0
9
In particular, fractal growth may be observed by a mechanism similar to this one.

Or, which amounts to the same thing, a metallic cylinder with infinite length, on which
we assume translation invariance.
10
171
Peter Gustav Lejeune-Dirichlet (18051859), German mathematician, was professor successively in Breslau, Berlin, and Gttingen, where he succeeded his teacher Gauss (see page 557). He
was also the brother-in-law of Jacobi and a friend of Fourier. He
studied trigonometric series and gave a condition for the convergence of a Fourier series at a given point (see Theorem 9.46 on
page 268). His other works concern the theory of functions in
analysis and number theory.
To prove this theorem, we start by stating a technical lemma which tells us,
in substance, anything that can be proved for the central point of the open
disc B(0 ; 1) can also be proved for an arbitrary point of the disc.
def
za
. Then a
1az
B = B(0 ; |a|1 ))
LEMMA 6.22 Let a B (0 ; 1). Denote a (z) =
B is some neighborhood of B(0 ; 1) (for instance
i) a is bijective on B ;
ii) a
B(0 ; 1)
H (B ), where
and satisfies
is a bijection from B(0 ; 1) into itself;
iii) a1 = a ;
iv) a (a) = 0.
In other words, a transforms the disc into itself, but moves the point a to the center in
a holomorphic manner.
Proof of Theorem 6.21. We use the mean value property for harmonic functions:

Z 2
Z

u a ( )
1
1
u(a) = u a (0) =
u a (e i ) d =
d .
2 0
2i B
With the change of variable = a ( ), we obtain = a ( ),

d
1 |a|2
= a ( ) =
,
d
(1 a )2
and
u(a) =
1
2i
u( )
1
a ( ) d =
a ( )
2i
u(e i )
1 |a|2
|a e i |
d.
In Theorem 6.21, put a = r e i ; there follows the equivalent formula:

Z 2

i
u(e i ) P (r, ) d,
u re =
0
where P is the Poisson kernel:
Conformal maps
172
Simon Denis Poisson (17811840) lost his father at fifteen, and

was sent by his mother to study medicine in Paris to earn his
living. Seeing patients die under his care, he decided (accidentally,
having first chosen the wrong classroom!) to dedicate himself to
mathematics. He was received first at the cole Polytechnique
at eighteen. Then honors came until the end of his life (he
was even made a baron by Louis XVIII). In mathematics, his
interests ranged over almost all of analysis: integration, Fourier
series, probability, and, especially, mathematical physics, celestial
mechanics and electrostatics in particular. A member of the Royal
Council of Public Education, he tried to develop the teaching of
mathematics.
DEFINITION 6.23 The Poisson kernel is the function P : [0, 1[R R
defined by
P (r, ) =
1 r2
1 1 r2
1
=

.
2 1 + r 2 2r cos
2 r e i 2
Remark 6.24 In the case of holomorphic functions, the Cauchy formula not only reproduces
holomorphic functions, but also creates them: if f is a given function on the circle C =
B(0 ; 1), assumed to be merely continuous, then
Z
1
f ( )
def
F (z) =
d
2i C z
is a holomorphic function on B(0 ; 1).

However, there isnt a very strong link between f and F , contrary to what one might
expect. For instance, for z0 C , nothing forces that f (z0 ) = limzz0 F (z). One may, for
instance, take f (z) = z (which, by the way, coincides on the unit circle with the function
g(z) = 1/z). The preceding formula leads to F (z) = 0 for any z B(0 ; 1), as the reader can
amuse herself to prove.
On the other hand, for harmonic functions, there is a formula to create a harmonic
function from boundary values, with a continuity property, using the Poisson kernel. This is
the object of the next theorem.
THEOREM 6.25 (Dirichlet problem for a disc) Let f : B (0 ; 1) R be a

continuous function defined on the unit circle. Put
Z
1 2
1 |z|2
f (e i )
if z B(0 ; 1),
d
z e i 2
u(z) = 2 0
f (z)
if z B(0 ; 1).
Then u is continuous on B(0 ; 1) and harmonic on B(0 ; 1).

If f is has discontinuities, then u is harmonic on B(0 ; 1) and continuous in each
point of the unit circle in which f happens to be continous.
Proof. The proof of this theorem is omitted. Still, we remark that it involves a
sequence of Dirac functions given by the Poisson kernel for [r 1]. After reading
173
1
0,5
0
0,5
1
0,5
1
0,5
0
0,5
0,5
1 1
Fig. 6.5 The temperature is equal to 1 on the lower half-circle, and to +1 on the upper
half-circle; it is a harmonic function on the disc. Here it is represented on the
vertical axis, and is deduced from the preceding theorem.
chapter 8 on distributions, and in particular section 8.2.b on page 231, the reader is
invited to try to formalize the proof of this theorem.
The Dirichlet problem has many applications in physics. Indeed, there are
many equations of harmonic type: electromagnetism, hydrodynamics, heat
theory, and so on. The example already mentioned find the temperature at
any point in a disc where the temperature is imposed to be T0 on the upper
half-circle and T1 on the lower half-circle is illustrated in Figure 6.5. After
integration, one finds (in polar coordinates)

2
2r sin
T (r, ) = arctan
.
1 r2
In Exercises 6.7 and 6.8, two other examples of solutions of the Dirichlet
problem are given, for a half-plane and for a strip, respectively.
Conformal maps
174
EXERCISES
Exercise 6.1 Let U : R2 R2 be a twice differentiable function. Show that if f : C C is
holomorphic, then denoting f : (x + i y) 7 (X + i Y ), we have

2
e (X , Y ),
(x, y) U (x, y) = f (x + i y) (X ,Y ) U
def
e (X , Y ) =
with U
U f 1 (X , Y ).
Exercise 6.2 Show that any two annuli are homeomorphic (one can use, for instance, a
purely radial function).
Complex flows
Exercise 6.3 Determine the flow of the wind in a two-dimensional domain delimited by two
half-lines meeting at an angle (for instance, one of the half-lines represents the horizontal
ground, and the other a mountainside with uniform slope). Use the following boundary
conditions: the wind must have velocity with modulus v and flow parallel to the boundaries,
infinitely far away from the intersection of the lines.
Assume that the flow is irrotational, incompressible, and stationary. In particular, what is
the velocity of the wind at the point where the slope changes? Consider the various possibilities
depending on the value of [0, ].
Exercise 6.4 Let k be a real number. Interpret the complex flows given by
1 (z) = k log(z a)
and
2 (z) = ik log(z a).
Exercise 6.5 We wish to solve the problem of the flow of a Newtonian fluid in three dimen-
sions around a cylinder with infinite length and circular section. If only the solutions which
are invariant by translation along the axis are sought, the problem is reduced to the study of
the flow of a fluid in two dimensions around a disc.
We assume that the velocity of the fluid, far from the obstacle, is uniform. We look for
the solution as a complex potential written in the form
= V0 z + G (z),
where the term V0 z is the solution unperturbed by the disc and where G (z) is a perturbation,
such that
G (z) 0.
|z|+
Check that the solution obtained is indeed compatible with the boundary conditions (the
velocity of the fluid is tangent to the disc), and find the stagnation points of the fluid.
Exercises
175
Dirichlet problem
Exercise 6.6 We wish to know the temperature of a sheet in the shape of an infinite half-
plane, represented in the complex plane by the closed region

D = {z C ; Im(z) 0}.
We assume that the temperature on the boundary of the sheet is given by
T0 if x < 1,
T (x) = T1 if |x| < 1,
T2 if x > 1.
If z D, we denote by 1 and 2 the angles defined in the picture below:

z
D
1
1
2
1
Show that any function z 7 1 + 2 + (with , , R) is harmonic. Deduce from

this the solution to the Dirichlet problem.
Exercise 6.7 Dirichlet problem for a half-plane: consider the domain
We want to find u of
to D and such that
C2
D = {z C ; Imz > 0}.
class defined on D with values in R, admitting a continuous extension
u = 0
in D,
u| D = u0 given.
Find, using the Cauchy formula, an integral formula for u(z) when Imz > 0.
Hint: Show, before concluding, that if z = x + i y with y > 0 and if f is holomorphic
on D, then
Z
y f ()
1 +
d.
f (z) =
( x)2 + y 2
Exercise 6.8 Dirichlet problem for a strip: we are seeking the solution of the Dirichlet
problem for a strip
B = {z C ; 0 < Imz < 1} .

In other words, we are looking for a function u of C 2 class defined on B with values in R,
admitting a continuous extension to B, and such that
u = 0
in B,
u(x) = u0 (x) x R,
u(x + i) = u1 (x) x R.
def
Show that if u = Re f (with f holomorphic) is a solution of the problem, then we have

Z

Z

i +
(t z)
i +
(t z)
f (z) =
u0 (t) coth
dt +
u1 (t) tanh
dt.
2
2
2
2
176
Conformal maps
SOLUTIONS
Solution of exercise 6.2. Writing as usual z = r e i and denoting by r1 < r < r2 the first
annulus and by r1 < r < r2 the second one, consider the map

r r1
r 7 2
(r r1 ) + r1 ,
r2 r1
7 .
This application is indeed an homeomorphism from the first annulus onto the second (the
transformation of the radius is affine). One can check that this is not a holomorphic map,
except if the two annuli are homothetic to each other, in which case the formula given is the
corresponding homothety.
Solution of exercise 6.3. We can transform the domain D into the upper half-plane by
means of the following conformal map:
z 7 w = z / .
In the domain D = {z C ; Im(z) > 0}, the flow of the wind is given by a complex potential
(w) = vw; one checks indeed that the velocity field that derives from this is uniform and
parallel to the ground (since v R). The function being holomorphic, its real part Re,
which is the real velocity potential, is harmonic: (Re) = 0. Coming back to the initial
domain, we put

(z) = w(z) = z / = v z / ,

which therefore satisfies (Re = 0 (because of Corollary 6.15) and has the required boundary conditions. The velocity of he wind is given by the gradient of Re.
At z = 0, the mapping z 7 w is not conformal. If < , the derivative
/1
dw
=
z
dz
tends to 0 (and so does the velocity of the wind), which shows that at the bottom of a Vshaped valley one is protected from the transversal wind (but not from the longitudinal wind,
as people from Lyons know, who are victim of the North-South wind, whereas the East-West
wind is almost absent). On the other hand, if > , this derivative tends to infinity, which
shows that on the crest of a mountain, the transverse wind is important.
Solution of exercise 6.4. 1 corresponds to the flow around a source or a sink (depending
on the sign of k), and 2 corresponds to a whirlpool, as one can see by picturing the velocity
field.
Solution of exercise 6.5. Rotating the entire system if necessary, we may assume that the
velocity at infinity is parallel to the real axis. Therefore we take V0 R and look for a flow
symmetric with respect to the real axis.
We know a conformal transformation that maps the half-plane minus a half-disc into the
upper half-plane; it is given by z = z + a 2 /z, where a is the radius of the disc that is removed.
The circle C (0 ; a) is mapped to the segment [2a, 2a].
There only remains to find a free flow in the new domain: it is given trivially by
(w) = V0 w.
The solution of the problem is therefore

a2
(z) = V0 z +
.
z
177
Notice that it is of the stated perturbed form, with G (z) = V0 a 2 /z.

One can show that, if |z| = a, then (z) is tangent to the circle (i.e., (z) is orthogonal
to the circle).
The stagnation points of the fluid are given by (z) = 0, which gives z = a.
Solution of exercise 6.6. Notice that the function which is proposed is the imaginary part
of
z 7 L(z + 1) + L(z 1) + ,
where L is the holomorphic logarithm defined on C with a cut along the lower half-plane. We
look for the constants , , and by noticing that on the boundary of the sheet, T = T0 for
1 = 2 = , T = T1 for 1 = 0 and 2 = , and T = T2 for 1 = 2 = 0. We obtain then
T T2
T0 T1
1 + 1
2 + T2
T T1
y
T T2
y
= 0
arctan
+ 1
arctan
+ T2 .
x +1
x 1
T =
Chapter
Distributions I
7.1
Physical approach
7.1.a
The problem of distribution of charge
We know that a point-like particle with electric charge q, placed at a point r 0

in the usual space R3 , produces, at any point r in R3 \ { r 0 }, an electrostatic
field
q
r r0
E( r ) =
.
()
40 k r r 0 k3
If more charges are present, the linearity of the Maxwell equations ensures
that the fields created by all the charges is the sum of those created by each
charge. However, when working in macroscopic scale, it is sometimes better to
describe the distribution of charges in continuous form; it is then modeled by
a function : R3 R which associates, to each point, the density of electric
charge at this point. The interpretation of the function is the following:
if r 0 R3 and if d3 r = dx dy dz is an elementary volume around r 0 , then
( r 0 ) dx dy dz
represents the total charge contained in the elementary volume dx dy dz.
Distributions I
180
The electric field at a point r is then (assuming the integral exists):

ZZZ
1
r r
E( r) =
( r )
d3 r ,
k3
40
3
k
r
r
R
()
and the total charge contained in a volume V R3 is

ZZZ
Q=
( r ) d3 r.
V
Can we reconcile the two expressions () and () for the point-like charge
and the density of charge respectively? In other words, can we write both
equations in a uniform way?
To express the point-like charge with the continuous viewpoint, we
would need a function r 0 describing a point-like unit charge located at r 0 .
This function would thus be zero on R3 \ { r 0 }; moreover, when integrated
over a volume V , we would obtain
ZZZ
1 if r 0 V ,
3
r0 ( r ) d r =
0 if r 0
/V.
V
Generalizing, for any continuous function f : R3 R, we would need to
have
ZZZ
f ( r ) r 0 ( r ) d3 r = f ( r 0 ).
(7.1)
V
But, according to the theory of integration, since r 0 is zero almost everywhere, the integral (7.1) is necessarily zero also. A Dirac function cannot
exist.
One may, following Dirac, use a sequence of functions with constant integral equal to 1, positive, and that concentrate around 0, as for instance:
n if |x| 1/2n,
n (x) =
0 if |x| > 1/2n,
or its equivalent in three dimensions
3n3 /4
Dn ( r ) =
0
if k rk 1/2n,
if k rk > 1/2n,
Every time the temptation arises to use the -function, the sequence of
functions (n )nN (resp. (Dn )nN ) can be used instead, and at the very end of
the computation one must take the limit [n ]. Thus, r 0 is replaced in (7.1)
by Dn ( r r 0 ), and the formula is written
ZZZ
ZZZ
3
f ( r ) r 0 ( r ) d r = lim
f ( r ) Dn ( r r 0 ) d3 r = f ( r 0 ).
V
Physical approach
181
The reader can check that this formula is valid for any function f continuous
at r 0 (see for instance a fairly complete discussion of this point of view in the
book by Cohen-Tannoudji et al. [20, appendix II]). However, this procedure
is fairly heavy. This is what motivated Laurent Schwartz (and, independently,
on the other side of the Iron Curtain, Isral Gelfand [38]) to create the
theory of distributions (resp. of generalized functions), which gives a rigorous meaning to the -function and justifies the preceding technique (see
Theorem 8.18 on page 232).
In the same manner we will want to describe not only the distribution of
point-like charges, but also a charge supported on a surface or a curve, so that,
if we denote by ( r ) the function giving the distribution of charges, we can
compute the total charge contained within the volume V with the formula
ZZZ
Q (V ) =
( r) d3 r
V
Z
ZZ
ZZZ
X
2
=
qi +
d +
d s+
vol. ( r) d3 r ,
charges i in
the volume V
L V
S V
where vol. is the volumic density of charge, the surface distribution on the
surface S , the linear distribution on the curve L , and q i the point-like
charges.
7.1.b The problem of momentum

and forces during an elastic shock
In this second example, we will look at the problem of an elastic shock between
two objects.
Consider a game of squash. Assuming that at a given time the ball reaches
a wall (orthogonally to the surface, to simplify) with velocity v0 , and rebounds.
The ball gets squashed a little bit, so that the contact lasts for a period of
time t which is nonzero; then it rebounds with velocity v0 . Hence the
graph of the speed of the ball as a function of time has the following shape:
v
t
Newtons law stipulates that, during the movement, the force f exerted on
the ball is such that f = m
v ; it is therefore proportional to the derivative
of the function graphed above. Now if we wish to model in a simple way a
hard collision (where the squashing of the ball is not taken into account,
for instance, as in a game of ptanque), the graph of the speed becomes
182
Distributions I
v
t
and the force exerted should still be proportional to the derivative of this
function; in a word, f should be zero for any t 6= 0 and satisfy
Z
1 +
f (t) dt = v(+) v() = 2v0 .
m
Once again, neither the integral nor the previous derivative can be treated in
the usual sense of functions.
7.2
Definitions and examples of distributions
We will now present a mathematical tool which, applied to a continuous function f , gives its value at 0, a relation which will be written
( f ) = f (0)
or
, f = f (0).
The notation , f is simply a convenient equivalent to ( f ) so this means

that is a machine that acts on functions, associating them with a number.
Now the rule above happens to be linear as a function of f ; so it is natural
to apply it to elements of a certain vector space of functions (that may still be
subject to some choice). This object is thus what is called a complex linear
form or linear functional on a space of functions.1 So we have
: space of functions C,
function 7 number.
There remains to make precise on which space of functions the functional
(and other useful functionals) will be defined.
DEFINITION 7.1 (Test function) Let n 1. The test space, denoted D(Rn ),
is the vector space of functions from Rn into C, which are of class C and
have bounded support (i.e., they vanish outside some ball; in the case n = 1,
they vanish outside a bounded interval).
A test function is any function D(Rn ).
1
Bourbaki presents the operation of integration as a linear form on a space of functions.

According to L. Schwartz, this abstract approach seems to have played a role in the genesis of
the theory of distributions.
183
When the value of n is clear or irrelevant in the context, we will simply

denote D instead of D(Rn ).
Example 7.2 There are an abundance of test functions, but it is not so easy to find one. (In
particular because they are not analytic on Rn can you see why?)
As an exercise, the reader will show that for a, b R with a < b , the function

1
exp
if x ]a, b [ ,
(x) =
(x a)(x b )
0
if x ], a] [b , +[
is a test function on R.
DEFINITION 7.3 A distribution on Rn is any continuous linear functional
defined on D(Rn ). The distributions form a vector space called the topological dual of D(Rn ), also called the space of distributions and denoted D (R
Rn)
or D .
For a distribution T D and a test function D, the value of T at
will usually be denoted not T () but instead T , . Thus, for any T D
and any D, T , is a complex number.
Remark 7.4 This definition, simple in appearance, deserves some comments.
Why so many constraints in the definition of the functions of D? The reason is simple:
the topological dual of D gets bigger, or richer, if D is small. The restriction of the
space test D to a very restricted subset of functions produces a space of distribution which is
very large.2 (This is therefore the opposite phenomenon from what might have been expected from
the experience in finite dimensions.3)
The C regularity is necessary but the condition of bounded support can be relaxed, and
there is a space of functions larger than D which is still small enough that most interesting
distributions belong to its dual. It is the Schwartz space S , which will be defined on page 289.
A distribution T is, by definition, a continuous linear functional on D; this means that for
any sequence (n )nN of test functions that converges to a test function D, the sequence
of complex numbers (T , n )nN must converge to the value T , . To make this precise we
must first specify precisely what is meant by a sequence of test functions (n )nN converging
to .
DEFINITION 7.5 (Convergence in D ) A sequence of test functions (n )nN

in D converges in D to a function D if
the supports of the functions n are contained in a fixed bounded

subset, independent of n;
2
For those still unconvinced:
we may define a linear functional on D, which to any test
R
function associates (t) dt. However, if we had taken as D a space which is too large, for
instance, the space L1loc of locally integrable functions (see Definition
7.8 on the next page), then
R
this functional would not be well-defined since the integral f (t) dt may have been divergent,
R
for instance, when f (t) 1. The functional 7 (t) dt is thus well-defined on the vector
space D of test functions, but not on the space of locally integrable functions. Hence it belongs
to D but not to (L1loc ) .
3
Recall that in finite dimensions the dual space is of the same dimension as the starting
space. Thus, the bigger the vector space, the bigger its dual.
Distributions I
184
all the partial derivatives of all order of the n converge uniformly to

the corresponding partial derivative of : for n = 1, for instance, this
means
( p) cv.u.
n ( p)
for any p N.
Now the continuity of a linear functional T defined on D can be defined

in the following natural manner:
DEFINITION 7.6 (Continuity of a linear functional) A linear functional T :
D C is continuous if and only if:
for any sequence ( )
n nN of test functions
in D,
if (n )nN converges in D to ,
the sequence T , converges in C to

n
T , .
Example 7.7 We can now show that the functional defined by : 7 (0) is continuous.
cv.u.
Let (n )nN be a sequence of test functions, converging to in D. Hence n , which

implies the simple convergence of (n )n to ; in particular, n (0) tends to (0) in C. This
shows that , n tends to , in C, which establishes the continuity of .
7.2.a
Regular distributions
Distributions can be seen as a generalization of the notion of functions or,

more precisely, of locally integrable functions.
DEFINITION 7.8 A measurable function f : Rn C is locally integrable if,
for any compact K Rn , the function f K is integrable on K . The space

of locally integrable functions on Rn is denoted L1loc (Rn ).
Functions which are locally integrable but not integrable are those which
have integrability problems at infinity.
Example 7.9 The function x 7 1/
|x| is locally integrable (but not integrable on R, for it

decreases too slowly). More generally, any function in L1 (Rn ), L2 (Rn ), or L p (Rn ) with p 1,
is locally integrable.
Any locally integrable function defines then a distribution by means of

the following theorem:
For any locally integrable function f ,
there is an associated distribution, also denoted f , defined by
Z
def
D
f , = f (x) (x) dx.
THEOREM 7.10 (Regular distributions)
185
Such a distribution is called the regular distribution associated to the locally integrable function f .
R
Proof. It suffices to show that the map 7 f (x) (x) dx is indeed linear (which
is obvious) and continuous.
Let (n )nN be a sequence of test functions converging to D. There exists a
bounded closed ball B such that B contains the supports
of all the n . Since the
R
function f is locally integrable, we can define M = B | f | and this M is a finite real
number.
R

R
We then have f (n ) M |n | M V kn k , where V is the
finite volume of the ball B. But kn k tends to 0 by the definition 7.5 and,
therefore, f , n tends to 0. By linearity, f , n tends to f , .
Remark 7.11 If the functions f and g are equal almost everywhere, the distributions f and g
will be equal, that is, f , = g, for any D. The converse is also true: if two regular
distributions are equal as functionals, then the corresponding functions are equal almost everywhere. The reader is invited to construct a proof of this fact using the following property:
any integrable function with compact support can be approximated in the mean by functions
in D. So the notion of distribution is in fact an extension of the notion of classes of locally
integrable functions equal almost everywhere.
Example 7.12 Denote by 1 the constant function 1 : x 7 1 (defined on R). The regular
distribution associated to 1 is thus the map
1 : D C,
7 1, =
(x) dx.
Regular distributions will have considerable importance throughout this

chapter, since they are those which look the most like classical functions.
They will be used more than once to discover a new definition when time will
come to generalize to distributions certain concepts associated to functions
(such as the derivative, the Fourier transform, etc.).
7.2.b Singular distributions

We now define the Dirac distribution:
DEFINITION 7.13 The Dirac distribution is the distribution which, to any
function in D(R), associates its value at 0:
def
, = (0)
D.
For a R, we define similarly the Dirac distribution centered at a, denoted a , by its action on any test function:
def
a , = (a)
D.
Frequently (this will be justified page 189), a will instead be denoted (x a).
The preceding definition is easily generalized to the case of many dimensions:
Distributions I
186
The English physicist and mathematician Paul Adrien Maurice

Dirac (19021984) was one of the great founders of quantum
mechanics. One of his major achievements was to give, for the
first time, a correct description of a quantum relativistic electron
(Dirac equation), interpreting the solutions with negative energy
that were embarrassing physicists as antiparticles. He also showed
how the mathematical theory of groups could bring precious
information concerning particle physics. According to Dirac, the
mathematical beauty of a physical theory should be of paramount
importance.
DEFINITION 7.14 (Dirac Distribution on R n ) Let r 0 Rn be a point in Rn .
The Dirac distribution at r 0 is the distribution defined by its action on any

test function in D(Rn ):
def
r 0 , = ( r 0 )
D.
As before, it will often be denoted ( r r 0 ).
We have thus introduced a notation that may at first lead to some

confusion: if T is a distribution, T ( r) means not the value of T at the point r,
which would make no sense,4 but the distribution T itself; this serves as reminder that
the variable of the function on which T acts is denoted r.
Remark 7.15
DEFINITION 7.16 The Dirac comb is the distribution X (pronounced sha5)
defined by its action on any test functions D(R):

def
X, =
that is,
X=
+
X
n=
or
+
X
(n),
n=
def
X(x) =
+
X
n=
(x n).
P
Notice that the distribution X is well defined, because the sum n (n)
is in fact finite ( being, by definition, with bounded support).
The continuity is left as an (easy) exercise for the reader.
The distribution is represented graphically by a vertical arrow, as here:
4
This bears repeating again at the risk of being boring: a distribution is not a function, but
a linear form on a space of functions.
5
This is a Cyrillic alphabet letter, corresponding to the sound sh.
187
a
x
a x
a 12 a + 12
The distribution
The X distribution
7.2.c Support of a distribution

Two distributions S and T are equal if S , = T , for any D. They
are equal on an open subset in R n if S , = T , for any D such
that the support of is contained in (i.e., for any D which is zero
outside ).
Example 7.17 Let H be the Heaviside function. The regular distributions 1 and H are equal
on ]0, +[. Indeed, for any function D which vanishes on R , we have
Z +
Z +
1, =
1 (x) dx =
(x) dx = H , .
Similarly, X and are equal on ] 21 , 12 [.
Consider now all the open subsets on which a given distribution T is

zero (i.e., is equal to the zero distribution on ). The union of these open sets
is itself an open set. One can then show (but this is not obvious and we will
omit it here, see [80] for instance) that this is the largest open set on which T
is zero. Its complement, which is therefore closed, is called the support of the
distribution T , and is denoted supp(T ).
Example 7.18 The support of the Dirac distribution centered at 0 is supp() = {0}, and more
generally, for a R, we have supp(a ) = {a}.

The support of the Dirac comb is supp(X) = Z.
7.2.d Other examples

PROPOSITION 7.19 Let D(R) be a test function. The quantity
has a finite limit as [
(x)
dx +
x
(x)
dx
x
Moreover, the map

Z
(x)
7 lim+
dx
0
|x|> x
0+ ].
(7.2)
is linear and continuous on D(R).

Proof. Because of its length, the proof of this theorem is given in Appendix D,
page 600.
Distributions I
188
DEFINITION 7.20 The Cauchy principal value of 1/ x is the distribution on
R defined by its action on any given by

Z
Z
(x)
(x)
1
def
pv , = pv
dx = lim+
dx.
x
x
0
|x|> x
It should be emphasized that this integral is defined symmetrically around the

origin. It is this symmetry which provides the existence of the limit defining
the distribution. The distribution pv 1/x can also be seen as the derivative
of the regular distribution associated to the locally integrable function x 7
log |x| (see Exercise 8.4 on page 241).
Other properties and examples of physical applications of the Cauchy
principal value will be given in the next chapter. In addition, a generalization
of the method leading to the definition of a distribution from the function
x 7 1/x (which is not locally integrable) is given in Exercises 8.5 and 8.6 on
page 241.
7.3
Elementary properties. Operations
7.3.a
Operations on distributions
We wish to define certain operations on distributions, such as translations,

scaling, derivation, and so on. For this purpose, the method will be consistently the same:
Consider how those operations are defined for a locally integrable function; write them, in the language of distributions, for the associated
regular distribution; then generalize this to an arbitrary distribution.
For instance, for any locally integrable function f on Rn and any a Rn ,
we define the translate fa of the function f by
x R
def
fa (x) = f (x a).
In other words, the graph (in Rn+1 ) of the function f has been translated by
the vector a:
189
fa
a
The regular distribution associated to this function, which is also denoted
fa , therefore satisfies
Z
Z

fa , = f (x a) (x) dx = f ( y) ( y + a) dy = f , a .
This can be generalized to an arbitrary distribution T in the following manner:
DEFINITION 7.21 The translate of the distribution T by a, denoted Ta or
T (x a), is defined by

def
Ta , = T , a
which is also written
for any D(Rn ),
def

T (x a), (x) = T (x), (x + a) .

The reader
can easily check that the number T , a exists and that the map
7 T , a is indeed linear and continuous.
Remark 7.22 As in Remark 7.15, the notation T (x a) is simply synonymous with Ta ; similarly, with an abuse of language, (x a) sometimes designates (for physicists more than for
mathematicians) the function t 7 (t a) and not the value of at x a.
Example 7.23 Let a R. The translate by a of the Cauchy principal value is
pv
1
: 7 lim+
0
x a
(x)
dx +
x a
a+
(x)
dx .
x a
By analogous reasoning, we define the transpose of a distribution. For

a function f : Rn C, we denote by f the transpose of f , which is the
def
function x 7 f(x) = f (x). If f is locally integrable, then so is f; hence
both f and f define regular distributions. Moreover, for any D(Rn ), we
have
Z
Z
Z
x=t
f , = f (x) (x) dx = f (x) (x) dx =

f (t) (t) dt = f ,
.
This justifies the following definition:
190
Distributions I
DEFINITION 7.24 (Transpose) The transpose of a distribution T on Rn , de-
noted T or T (x), is defined by
,
T , = T ,
def
which can also be written:
def

T (x), (x) = T (x), (x)
for any D(Rn ).
Similarly, for any a R , a change of scale defines the dilation by a of a

function f by x 7 f (a x). The regular distribution associated to the dilation
of a function satisfies
Z
Z
x dx

f (a x), (x) = f (a x) (x) dx = f (x)
a |a|
x E
1 D
=
f (x),
.
|a|
a
This leads to the following generalization:
DEFINITION 7.25 (Dilation) Let a R . The dilation of a distribution T by
the factor a, denoted T (a x), is defined by

x E
def 1 D
T (a x), (x) =
T (x),
|a|
a
for any D(R).
Dilations of a distribution on Rn are defined in exactly the same manner.

However, the jacobian associated to the change of variable y = a x is |a|n
instead of |a|, which gives the dilation relations
x E

1 D
T (a x), ( x) = n T ( x),
on Rn .
|a|
a
There is, unfortunately, no general definition of the product of two distributions.6 Already, if f and g are two locally integrable functions (hence
defining regular distributions), their product is not necessarily locally integrable.7 However, if is a function of class C , then for any D, the
product is a test function (still being C and with bounded support); if
f is locally integrable, then f is also locally integrable. From the point of
view of distributions, we can write
Z
Z

f , =
(x) f (x) (x) dx = f (x) (x) (x) dx = f , .
6
In particular, it is easy to notice that taking the square of the Dirac distribution raises
arduous questions. However, it turns out to be possible to define a product in a larger space
than the space of distributions, the space of generalized distributions of Colombeau [22, 23].
p
7
Simply take f (x) = g(x) = 1/ x.
191
This leads us to state the following definition:

DEFINITION 7.26 (Product of a distribution by a C function) Let T D
be a distribution and let be a function of class C . The product T D

is defined by
def
T , = T ,
for any D.
By applying this definition to the Dirac distribution, the following fundamental result follows:
THEOREM 7.27 Let be a function of C class. Then we have
(x) (x) = (0) (x)

Similarly, we have the relation
(x) X(x) =
and in particular
X
nZ
x (x) = 0.
(n) (x n).
Proof. If is of C class, then by definition

, = , = (0) (0) = (0) , = (0) , .
Since any test function has bounded support, the last formula follows by simple linearity.
It is important to notice that, for instance, the product H makes no

sense here!
THEOREM 7.28 The equation x T (x) = 0, with unknown a distribution T , admits
the multiples of the Dirac distribution as solutions, and they are the only solutions:

x T (x) = 0 T = , with C .
Remark 7.29 If is merely continuous, the product T cannot be defined for an arbitrary
distribution T , but the product can in fact still be defined. Similarly, if is of C 1 class
(in a neighborhood of 0), the product can be defined (see the next section).
7.3.b Derivative of a distribution

The next goal is to define derivation in the space of distributions. As before,
consider a locally integrable function f L1loc , and assume moreover that it
is differentiable and its derivative is also locally integrable. Then the regular
distribution associated to f is given by
Z
Z
f , = f (x) (x) dx = f (x) (x) dx = f , .
The passage to the second integral is done by integrating by parts; the boundary terms vanish because has bounded support (this is one reason for the
very restrictive choice of the space of test functions). This justifies the
following generalization:
Distributions I
192
DEFINITION 7.30 The derivative of a distribution T D (R

R) is defined by
def

T , = T ,
for any D(R).
(7.3)
Similarly, the higher order derivatives are defined by induction by

def
T (m) , = (1)m T , (m)

for any D and any m N.
Example 7.31 The function x 7 log |x| is locally integrable and therefore defines a regular
distribution. However, its derivative (in the usual sense of functions) x 7 1/x is not locally
integrable; therefore it does not define a (regular) distribution.

However, if the derivation is performed directly in the sense of distributions, log |x| is a
distribution (it is shown in Exercise 8.4 on page 241 that it is the Cauchy principal value pv 1x ).
If working on Rn , a similar computation justifies the definition of partial

derivatives of a distribution T D (Rn ) by

T
def
, = T,
xi
xi
for D(Rn ) and 1 i n. Partial derivatives of higher order are obtained
by successive applications of these rules: for instance,
3

T
3
3
, = (1) T , 2
.
x12 x2
x1 x2
In particular, on R3 , we have

T
, = T ,
,
, = T ,
,
x
x
y
y

T
, = T ,
.
z
z
Similarly, the laplacian being a differential operator of order 2, we define:
DEFINITION 7.32 The laplacian of a distribution on Rn is given by
def
T , = T ,
for any D(Rn ).
THEOREM 7.33 (Differentiability of distributions) Any distribution on Rn ad-
mits partial derivatives of arbitrary order. Any partial derivative of a distribution is

a distribution.
COROLLARY 7.33.1 Any locally integrable function on R defining a distribution,
it is infinitely differentiable in the sense of distributions:

f L1loc (R)
k N
f (k) D (R).
Dirac and its derivatives
193
At the age of sixteen, Oliver Heaviside (18501925), born in a

very poor family, had to leave school. At eighteen, he was hired
by Charles Wheatstone, the inventor of the telegraph. He then
studied the differential equations which govern electric signals
and in particular the
p
p square root of the derivative (he wrote
d/dt H (t) = 1/ t), which gave apoplexies to the mathematicians of Cambridge.9 One of his papers occasioned the severe disapproval of the Royal Society. He defended vector analysis
(Maxwell was in favor of the use of quaternions instead of vectors), symbolic calculus, and the use of divergent series, and had
faith in methods which work before being rigorously justified.
When Laurent Schwartz presented the theory of distributions, this result

seemed particularly marvelous to the audience: all distributions are infinitely
differentiable and their successive derivatives are themselves distributions.8
Thus, if we consider a locally integrable function, however irregular, its associated distribution is still infinitely differentiable. This is one of the key
advantages of the theory of distributions.
7.4
7.4.a
The Heaviside distribution
A particularly important example is the Heaviside distribution, which we will

use constantly.
DEFINITION 7.34 The Heaviside function is the function H : R R de-
fined by
0
H (x) = 1/2
if x < 0,
if x = 0,
if x > 0.
8
Mathematicians had been aware for less than a century of the fact that a continuous function could be very far from differentiable (since there exist, for instance, functions continuous
on all of R which are nowhere differentiable, as shown by Peano, Weierstrass, and others).
9
Enamored of rigor from a recent date only and, therefore, as inflexible as any new convert.
The reader interested in the battle between Oliver Heaviside and the rigorists of Cambridge
can read with profit the article [48].
Distributions I
194
DEFINITION 7.35 The Heaviside distribution is the regular distribution as-
sociated to the Heaviside function. It is thus defined by

Z +
def
(x) dx.
D(R)
H , =
0
The derivative, in the sense of distributions, of H is then by definition

Z +
+
H , = H, =
(x) dx = 0 = (0),
0
from which we deduce
THEOREM 7.36 The derivative of the Heaviside distribution is the Dirac distribution
H =
in the sense of distributions.
7.4.b Multidimensional Dirac distributions

In many cases, the physicist is working in R3 or R4 . We will define here
the notions of point-like, linear, and surface Dirac distributions, which occur
frequently in electromagnetism, for instance.
Point-like Dirac distribution
Recall that in three dimensions the Dirac distribution is defined by
(3) def
, = (0).
The common notation (3) ( r ) = (x) ( y) (z) will be explained in Section 7.6.b. Most authors use instead of (3) since, in general, there is no
possible confusion. Thus a point-like charge in electrostatics is represented by
a distribution ( r ) = q ( r a). Note that, when it acts on a space variable,
the three-dimensional Dirac distribution has the dimension of the inverse of a
volume:
1
[] =
,
[L]3
which implies that ( r ) has the dimension Q /L3 of a volum density of
charge.
Surface Dirac distribution
To describe a surface carrying a uniform charge, we use the surface Dirac
distribution:
195
DEFINITION 7.37 Let S be a smooth surface in R3 . The (normalized) Dirac
surface distribution on S is given by its action on any test function

D(R3 ):
ZZ
def
S , =
d2 s,
S
where d2 s is the integration element over the surface S . In electromagnetism,

a uniform surface density 0 on S is described by the distribution 0 S .
This can be extended to nonuniform surface distribution, since we know how
to multiply a distribution by a function of C class. So a distribution of
charge ( r) of C class on a surface S is represented by ( r) S .
Example 7.38 In R3 , for R > 0, we denote by k xk R the Dirac surface distribution S ,
where S is the sphere S = x R3 ; k xk = R .
It should be noted that when working with spacial variables in R3 , the

surface Dirac distribution has the dimension of the inverse of length:
[S ] =
1
,
[L]
which can be recovered by writing (in a heuristic manner, indicated by large

quotes)
ZZZ
ZZ
3
S , =
S ( r ) ( r) d r =
( r) d2 s.
and let a R. Denote by S the surface obtained

from S by homothety with coefficient 1/a. Show that
1
S (a r) =
( r).
|a| S
Exercise 7.1 Let S be a surface in
R3
Solution: Let D(R3 ). Then we have

ZZ
r E
r
def 1 D
1
= 3
d2 r
S (a r), = 3 S ,
a
a
|a|
|a|
S
ZZ
1
1
= 32
( x) d2 x =
S , .
|a|
|a|
S
(definition)
Curvilinear Dirac distribution

By the same method, a density of charge on a curve is defined using the
curvilinear Dirac distribution:
DEFINITION 7.39 Let L be a curve in R3 .
The unit curvilinear Dirac

distribution on L is defined by its action on any test function D(R3 )
given by
Z
def
L , =
d,
L
196
Distributions I
+/k
k/2
k/2
/k
Fig. 7.1 The distribution can be seen as two Dirac peaks with opposite signs infinitely
close to each other that is, a dipole.
where d is the integration element over L . In electromagnetism, a density

of charge c ( r ) supported on the curve L is represented by the distribution c ( r ) S .
When working with spacial variables in R3 , the curvilinear Dirac distribution has the dimension of the inverse of a surface:
[L ] =
1
.
[L]2
7.4.c The distribution

As seen in the defining formula (7.3), the derivative of the Dirac distribution is defined by
def
, = (0)
for any test function D.
Note that if is a test function, then

, = (0) = lim

k
2

k2
k0
k

1
k
k
= lim
x
x +
, ,
k0 k
2
2
from which we conclude that (note the signs!)

1
k
k
x+
x
.
(x) = lim
k0 k
2
2
This formula is interpreted as follows. The distribution represents a positive
charge 1/k and a negative charge 1/k, situated at a distance equal to k, in
the limit where [k 0] (see Figure 7.1). Hence the distribution represents
a dipole, aligned on the horizontal axis (O x) with dipole moment 1, hence
oriented toward negative values of x.
197
Remark 7.40 This expression requires a notion of convergence in the space D , which will be
introduced in Section 8.2.a on page 230.
Can this be generalized to the three-dimensional case? In that situation, we

cannot simply take the derivative, but we must indicate the direction in which
to differentiate.
DEFINITION 7.41 In R3 , the distributions x , y , and z , are defined by
def
x , = x (0) = (0),
x
y , = y (0) = (0),
y
z , = z (0) = (0).
z
We denote by the vector ( x , y , z ).
Lets use these results to compute the electrostatic potential created by an elecrostatic
dipole. The potential created at a point s by a regular distribution of charges
is
ZZZ
1
( r )
V ( s) =
d3 r ,
40
k
r
sk
3
R
which can be written

1
1
.
V ( s) =
( r ),
40
k r sk
Now we extend this relation to other types of distributions, such as point-like,
curvilinear, surface, or dipolar distributions.
PROPOSITION 7.42 The potential created at a point s by the distribution of
charges ( r ) is
1
V ( s) =
40

( r ),
1
k r sk
(7.4)
Remark 7.43 The function s 7 1/ k r sk is not of C class, nor with bounded support.
As far as the support is concerned, we can work around the problem by imposing a
sufficiently fast decay at infinity of the distribution of charges. For the problem of the
singularity of 1/ k r sk at s = r, we note that the function still remains integrable since the
volume element is r 2 sin dr d d; this is enough for the expression (7.4) to make sense in
the case of most distributions of physical interest.
Consider now the dipolar distribution

T = P = P x x P y y Pz z ,
Distributions I
198
where we have denoted by P the dipole moment of the source. Then the
potential is given by

1
V ( s) = T ( r),
k r sk

1
= P x x P y y Pz z ,
k r sk

1
+ Py
+ Pz
= Px
,
x k r sk
y k r sk
z k r sk r=0
namely,
V ( s) = P x
sx
3
+ Py
sy
3
+ Pz
sz
3
Ps
k sk
k sk
k sk
k sk3
which is consistent with the result expected [50, p. 138].
THEOREM 7.44 In the usual space R3 , an electrostatic dipole with dipole moment
equal to P is described by the distribution P .
Exercise 7.2 Show that, in dimension one, dilating by a factor a 6= 0 yields
(a x) =
1
(x)
a |a|
(Solution page 245)
7.4.d Composition of with a function

The goal of this section is to give a meaning to the distribution f (x) ,
where f is a sufficiently regular function. This amounts to a change of
variable in a distribution and is therefore a generalization of the notions of
translation, dilation, and transposition.
Consider first the special case of a regular distribution and a change of
variable f : R R which is differentiable and bijective.
Let g be a locally integrable function. The function g f being still locally
integrable, it defines a distribution and we have
Z
Z

dy
g f , = g f (x) (x) dx = g( y) f 1 ( y) 1 ,
f f ( y)
since, in the substitution y = f (x), the absolute value of the jacobian is equal
to

dy 1
= f (x) = f f ( y) .
dx
By analogy, for a distribution T , we are led to define the distribution T f
as follows:
T
def
1
, f
T f , =
.
f f 1
199

Thus, the distribution f (x) satisfies

f 1 (0)
1
, f
f (x) , =
= 1 ,
f f 1
f f (0)
for any D, which can also be written

1
( y y0 ),
f (x) =
f ( y0 )
where y0 is the unique real number such that f ( y0 ) = 0. Note that it is

necessary that f ( y0 ) 6= 0: the function f must not vanish at the same time
as its derivative.
We can then generalize to the case of a function f not necessarily bijective
but still differentiable by writing (within quotes!)
Z +
X

1

(x).
f , =
f (x) (x) dx =
(x)

f
x | f (x)=0
To summarize:
THEOREM 7.45 Let f be a differentiable function which has only isolated zeros.
Denote by Z( f ) the set of these zeros: Z( f ) = { y R ; f ( y) = 0}. Assume that

the derivative f of f does not vanish at any of the points in Z( f ). Then we have
X

1

(x y).
f (x) =
f ( y)
yZ ( f )
A classical application of this result to special relativity is proposed in Exercise 8.9 on page 242.
7.4.e
Charge and current densities
In this section, we study the application of the preceding change of variable

formula, in a simple physical situation: the Lorentz transformation of a current four-vector.
In electromagnetism, the state of the system of charges is described by data
consisting of
the charge density ( x, t) 7 ( x, t) ;
the current density ( x, t) 7 j( x, t).
In the case of a point-like particle with charge q, if we denote by t 7 R(t) its

position and by t 7 V (t) its velocity, with V (t) = dR(t)/dt, these densities
are given by

( x, t) = q (3) x R(t)
and
j ( x, t) = q V (t) (3) x R(t) .
Distributions I
200
Moreover, the densities of charge and current form, from the point of
view of special relativity, a four-vector. In order to distinguish them, we will
denote three-dimensional vectors in the form j and four-vectors in the form
j = (, j /c ).
During a change of Galilean reference frame, characterized by a velocity
v (or, in nondimensional form, = v/c ), these quantities are transformed
according to the rule
( , j /c ) = (
) (, j/c ),
i.e., j = (
) j with
def
j = (, j /c ), (7.5)
where (
) is the matrix characterizing a Lorentz transformation. To simplify,
we consider the case of a Lorentz transformation along the axis (O x), characterized by the velocity = e x . Denote by x = (c t, x) = (x 0 , x 2 , x 2 , x 3 )
the coordinates in the original reference frame R and by x = (c t , x ) =
(x 0 , x 1 , x 2 , x 3 ) the coordinates in the reference frame R with velocity
relative to R. The Lorentz transformation of the coordinates is then given by

ct
ct

x
1
x
def
=
,
with = p
,
y
1
1 2
z
z
1
which we will write10 as x = (
) x. Expanding, we obtain

t = t x/c = (t v x/c 2 ),
t = (t + v x /c 2 ),
x = c t + x = (x v t),
x = (x + v t ),
and
y = y,
y= y,
z = z,
z = z.
Although it is almost always written in the form (7.5), the transformation

law of the current four-vector is more precisely given by

j(x ) = (
) j (
)1 x = (
) j(x),
or, in the case under consideration

( x , t) = ( x, t) 1c j x ( x, t) ,

j x ( x , t) = j x ( x, t) c ( x, t) .
(7.6)
The case of a particle at rest

Consider a simple case to begin with, where the particle is at rest in the
laboratory reference frame R.
10
Of course, tensorial notation can also be used to write x = x .
Derivation of a discontinuous function
201

Then j(x) = q (3) ( x), 0 . Looking back at equation (7.6), and denoting
by L 1 ( x ) = (x + v t ) the spacial part of the inverse transformation
)1 , we obtain
(

( x , t), 1c j ( x , t ) = q (3) L 1 ( x) ,
(3) L 1 ( x ) .
Notice then that

(3) L 1 ( x ) = (3) (x + v t ), y , z = (x + v t ) ( y ) (z )
= (x + v t ) ( y ) (z )
(change of variable)
= (3) ( x + vt)
which gives, as might be expected

( x , t ) = q (3) ( x + vt),
j ( x , t) = q v (3) ( x + vt).
So we see that the Dirac distribution has some kind of invariance property
with respect to Lorentz transformations: it does not acquire a factor ,
contrary to what equation (7.6) could suggest. This is a happy fact, since the
total charge in space is given by
ZZZ
( x , t) d3 x = q,
Q =
R3
and not q!
This result generalizes of course to the case of a particle with an arbitrary
motion compatible with special relativity.
7.5
7.5.a
Derivation of a function discontinuous at a point
We saw that any locally integrable function is differentiable in the sense of

distributions. Hence, we can differentiate the Heaviside function and find
that H = . If we had taken the derivative of H in the usual sense of
functions, we would have obtained a function zero everywhere except at 0,
where it is not defined. A primitive of this derivative would then have given
the zero function, not H .
In other words, the following diagram cannot be closed:
derivation
function H
y
derivation
y
no!
distribution H
202
Distributions I
The same will happen for any function with an isolated discontinuity at a
point: its derivative will exist in the sense of distributions, which differs from
the usual derivative (if it exists) by a Dirac peak with height equal to the jump
of the function.
The following notation will be used:
f is the function being studied, or its associated distribution;

f will be the regular distribution associated to the usual derivative
of f in the sense of functions:
Z +

f (x) (x) dx ;
f , =
this is of course a distribution; the derivative of f in the sense of

functions is a function which is not defined everywhere (in particular,
it is not at the points of discontinuity of f ), but the associated regular
distribution is well-defined as long as f is defined almost everywhere (and
locally integrable);
f will be the derivative of the distribution f , defined as above by
f , = f , ;
this is also a distribution.
Example 7.46 With this notation, the derivative in the sense of distributions of the Heaviside
distribution H is H = , but the distribution associated to the derivative of H taken in the

sense of functions is H = 0.
Let f be a function which is piecewise of C 1 class. Denote by a i the

(0)
points of discontinuity of f (we assume there are finitely many) and by i
(0)
the jump of f at a i : i = f (a +
i ) f (a i ). One can then write f as the sum
of a continuous function g and multiples of Heaviside functions (see Figure 7.2
on the next page):
X (0)
i H (x a i ).
f (x) = g(x) +
i
Since the usualderivatives

of the functions f and g are equal almost every
where, we have f = { g }. However, remembering that H = , we obtain,
in the sense of distributions
THEOREM 7.47 (Derivative of a discontinuous function) Let f be a function
which is piecewise C 1 . Then, with the notation above, we have

X (0)
f = f +
i (x a i ).
i
(7.7)
203
(0)
a
Fig. 7.2 Example of a function with a discontinuity at a with jump equal to (0) .
To lighten the notation, we will henceforth write equation (7.7) more compactly as follows:

f = f + (0) ;
all the discontinuities of f are implicitly taken into account. Similarly, if we
denote by (1) any discontinuity of the derivative of the function f , (2)
any discontinuity of its second derivative, and so on, then we will have the
following result, with similar conventions:
THEOREM 7.48 (Successive derivatives of a piecewise C function) Let f be
a piecewise C function. Then we have

f = f + (1) + (0) ,
and more generally
f (m) = f (m) + (m1) + + (0) (m1)
m N.
Example 7.49 Consider the function E : x 7 E (x) = integral part of x. Then
whereas, in the sense of distributions, we have
= X.

E = 0,
Exercise 7.3 Compute the successive derivatives of the function f : x 7 |x|.
Solution: Since f is continuous, its derivative in the sense of distributions

is the distribution associated to its usual derivative. But the derivative of f is
the function x 7 sgn(x) for any x 6= 0, which is defined
everywhere except at
x = 0. The associated regular distribution, denoted f , is the distribution
sgn(x) (meaning sign of x) defined by
Z +

sgn(x), (x) =
(x) (x) dx.
0

Therefore, the derivative of f in the sense of distributions is f = f = sgn.
The function x 7 sgn(x) has a discontinuity at 0 with jump equal to 2.
Moreover, its usual derivative is zero almost everywhere (it is undefined at 0
and zero everywhere else), which shows that
{ sgn } = 0
and
f = sgn = 2.
We deduce that the next derivatives of f are given by
f (k) = 2 (k2)
for any k 2.
Distributions I
204
7.5.b
Derivative of a function with discontinuity

along a surface S
Let S be a smooth surface in R3 , and let f : R3 C be a function which is

piecewise of C 1 class, with a discontinuity along the surface S .
We denote by x1 = x, x2 = y, and x3 = z the coordinates and assume the
surface S is oriented. For i = 1, 2, 3, we denote by i the angle between the
exterior normal n to the surface S and the axis (O x i ). Arguing as in the
previous section, the following theorem will follow:
THEOREM 7.50 (Derivatives of a function discontinuous along a surface) Let
f be a function with a discontinuity along the surface S . Denoting by (0) the jump
function S , equal to the value of f just outside minus the value just inside, we have

f
f
+ (0) cos i S .
=
(7.8)
xi
xi
n
i
xi
S
Proof. Assume that f is of class C 1 on R3 \ S and that its first derivatives admit
a limit on each side of S . We now evaluate the action of f / x on an arbitrary test
function D(R3 ):

ZZZ
f
, = f ,
=
f
dx dy dz
x
x
x

ZZ Z
=
f (x, y, z)
dx dy dz
x
for a smooth surface. Fix now y and z and denote by x the real number such that
(x , y , z ) S (generalizing to the case where more than one real number satisfies
this property is straightforward). In addition, put h(x) = f (x, y , z ). Then we have
Z
Z

h
f (x, y , z )
dx = h(x)
dx = h,
=
, ,
x
x
x
x
and we can now use Theorem 7.47, which gives

h
h
, =
+ (0) (x , y , z ) (x x ),
x
x
Z
h
=
dx + (0) (x , y , z ) (x , y , z ),
x
where we have denoted by (0) (x , y , z ) the jump of f across the surface S at the
point (x , y , z ).
205
But, integrating now with respect to y and z , we remark that

ZZ
ZZ
(0)
dy dz =
(0) cos x d2 s,
S
since d2 s cos x = dy dy, as can be seen in the figure.

z
dx dy
x
d2 s
y
x
To conclude, we write
ZZ

(0) cos x d2 s = (0) cos x S , .
Putting these results together, we obtain

ZZZ

f
f
, = (0) cos x S , +
dx dy dz
x
x
for any D(R3 ). The same argument performed with the partial derivatives with
respect to y or z leads to the stated formula (7.8).
Now we apply this result to a function f differentiable in a volume V

delimitied by a closed surface S = V , and equal to zero outside this volume.
Taking the test function (x) 1, we get11
ZZZ
ZZ
f
dv
f cos i ds = 0
for i = 1, 2, 3,
V xi
S
or, combining the three formulas for i = 1, 2, 3 in vector form, and noticing
that n = (cos x , cos y , cos z )
ZZZ
grad f d3 v =
V
ZZ
f n d2 s
Green-Ostrogradski formula.
11
More precisely, we should take a test function equal to 1 on V , but with bounded
support. If V is bounded, such a function always exists.
Distributions I
206
The formula (7.8) may also be applied to the components of a vector

function f = ( f1 , f2 , f3 ); then, after summing over indices, we get
ZZZ X
ZZ X
fi 3
d v
f i cos i d2 s = 0,
V i xi
S i
or equivalently
ZZZ
ZZ
div f d3 v =
( f n) d2 s
V
Green-Ostrogradski formula.
7.5.c Laplacian of a function discontinuous

along a surface S
By applying the results of the previous section, it is possible to compute, in
the sense of distributions, the laplacian of a function which is discontinuous
along a smooth surface.
Let S be a smooth
surface and f a function of C class except on the surface S . Denote by (0) the
jump of f across the surface S . Denote moreover by nf the derivative of f in the
normal direction n, defined by
f def X f
=
cos i = (grad f ) n,
n
xi
i
THEOREM 7.51 (Laplacian of a discontinuous function)
(1)
and let n be the jump of this function. Then we have

(1)
f = (0) n + n S + { f } ,
P
where n is the normal derivative of : n = i xi cos i = n.
(7.9)
Remark 7.52 The action of the distribution n on any test function is therefore given by

.
n , = S ,
n
THEOREM 7.53 (Green formula) Let f be a function twice differentiable in a
volume V bounded by a surface S , and let be a test function. Denoting by n the

exterior normal to S , we have

ZZZ
ZZ
f
3
(f f )d v =
f
d2 s.
(7.10)
n
V
S
Proof. Extend f by putting f ( r ) = 0 for any point outside V .

Note that the jump of f accross the surface S at a point r is f ( r ) when crossing
in the exterior direction, and similarly for the other jump functions involved.
207
Then apply directly the formula (7.9), noticing that

ZZZ
ZZ
(0)
2
d s,
f , =
f d3 v,
n, =
f
n
V
S
ZZZ
ZZ
(1)

f 2
d s,
{ f }, =
f d3 v.
n S , =
V
S n
7.5.d
Application: laplacian of 1/ r in 3-space
We continue working in R3 . If f is a function on a subset of R3 which is

twice differentiable and radial, its laplacian in the sense of functions can be
computed by the formula
1 2
(r f ).
r r2

Hence, in the sense of functions, we have 1r = 0, except possibly at r = 0.
We now try to compute this laplacian in the sense of distributions.
For any > 0, define
1/r for r > ,

f (r) =
0
for r < .
f =
Then apply the formula

(1)
f = (0) n + n S + { f }
to the function f . For this, denote by S the sphere r = and by V the

volume interior to this surface. The jumps of f and f / n accross the
surface S oriented toward the exterior are therefore equal to 1/r and 1/r 2 ,
respectively. We thus find
ZZZ
f , =
f d3 r
(by definition)
(1)
= (0) n + n S + { f } ,
but { f } = 0

ZZZ
ZZ
1
1
d2 s +
2 ( r) d3 r .
=
n
S
V
Each term in the last expression has a limit as [ 0], respectively, 0 (since
the derivatives of are bounded and the surface of integration has area of
order of magnitude 2 ) and
ZZZ
1
lim
2 ( r) d3 r = 4 (0)
0+
V
by continuity of (the factor 4 arises as the surface area of S divided
by 2 ).
Distributions I
208
Notice now that
ZZZ
1
1
, =
( r) d3 r
r
r
3
R
ZZZ
= lim+
f ( r ) ( r) d3 r = lim+ f , ,
0
R3
since the integrand on the left is an integrable function (because the volume
element is given by r 2 sin dr d d). Putting everything together, we find
that
ZZZ
1
1
, =
( r) d3 r = 4 (0).
r
R3 r
In other words, we have proved the following theorem:
THEOREM 7.54 The laplacian of the radial function f : r 7
= 4 .
r
1
r
1
k rk
is given by
It is this equation which, applied to electrostatics, will give us the Poisson

law (see Section 7.6.g on page 216).
In the same manner, one can show that
THEOREM 7.55 In R2 , the laplacian of r 7 log k rk is

log |r| = 2 .
Similarly, one gets (see [81]) :
PROPOSITION 7.56 For n 3, we have in Rn the formula
1
r n2
= (n 2)Sn ,
where Sn is the area of the unit sphere in Rn , namely

Sn =
(2)n/2
.
(n/2)
(7.11)
Convolution
209
7.6
Convolution
7.6.a
The tensor product of two functions
DEFINITION 7.57 (Tensor product of functions) Let f and g be two functions defined on R. The direct product of f and g (also called the tensor
product) is the function h : R2 R defined by h(x, y) = f (x) g( y) for any
x, y R. It will be denoted h = f g.
This definition is generalized in the obvious manner to the product of
f : R p C by g : Rn C.
A somewhat grandiloquent word (tensor product!!) for a very simple thing,

in the end. But despite the apparent triviality of this notion, the tensor
product of functions will be very useful.
Example 7.58 Let (x) = 1 for |x| < 1/2 and (x) = 0 elsewhere. The tensor product of
the functions H and is given by the following graphical representation, where the gray area
corresponds to the points where H takes the value 1, and the white area to those where it
takes the value 0:
y
1
2
H =
or
H (x, y) = H (x) ( y).
12
7.6.b
The tensor product of distributions
Lets now see how to generalize the tensor product to distributions. For this,
as usual, we consider the distribution associated to the tensor product of two
locally integrable functions. If f : R p C and g : Rn C are locally
integrable, and if we write h : R p+n C as their tensor product, then for any
test function D(R p+n ) we have
ZZ

h, = f (x) g( y), (x, y) =
f (x) g( y) (x, y) dx dy
Z
Z

D
E
= f (x)
g( y) (x, y) dy dx = f (x), g( y), (x, y) .
Take now a distribution S (x) on R p and a distribution T ( y) on Rn . The
function

def
x 7 (x) = T ( y), (x, y)
Distributions I
210
is a test function on Rn ; we can therefore let S act on it and, as for regular

distributions, we define the distribution S (x) T ( y) by
def D
E
S (x) T ( y), (x, y) = S (x), T ( y), (x, y) .
DEFINITION 7.59 (Tensor product of distributions) Let S and T be two distributions. The direct product, or tensor product, of the distributions S
and T is the distribution S (x) T ( y) defined on the space of test functions on
R p Rn by
def D
E
S (x) T ( y), (x, y) = S (x), T ( y), (x, y) .
It is denoted S T , or S (x) T ( y).
The tensor product is of course not commutative: S T 6= T S .

Remark 7.60 Be carefull not to mistake S (x) T ( y) with a product of distributions, which is
not well-defined in general (see note 6 on page 190).
Example 7.61 The Dirac distribution in two dimensions can be defined as the direct product
of two Dirac distributions in one dimension:

(2) ( r ) = (2) (x, y) = (x) ( y).
This can be illustrated graphically as follows:
( y)
y
(x)
(2) (x, y)
x
Similarly, in three dimensions, we have (3) ( x) = [ ](x, y, z) = (x) ( y) (z).
Example 7.62 In the same way, the tensor product H can be expressed by the graph
H ( y)
y
(x)
(x) H ( y)
Convolution
211
This distribution on R2 acts on a test function D(R2 ) by

Z +

(x) H ( y), (x, y) =
(0, y) dy.
0
Caution is that sometimes the constant function is omitted from the notation of a tensor product f 1. For instance, in two dimensions, the distribution denoted (x) is in reality equal to (x) 1( y) and acts therefore by the
formula
Z +

(x), (x, y) = (x) 1( y), (x, y) =
(0, y) dy.
This should not be confused with ( r ) = (x, y) = (x) ( y).
Example 7.63 (Special relativity) In the setting of Section 7.4.e, the density of charge (x) =
q (3) ( x) can also be written
since
= q 1 (3) ,

1 (3) (c t, x) = 1(c t) (3) ( x) = (3) ( x).
7.6.c
Convolution of two functions
DEFINITION 7.64 (Convolution of functions) Let f and g be two locally inte-
grable functions. Their convolution, or convolution product, is the function h

defined by
Z
def
f (t) g(x t) dt,
h(x) =
where this is well-defined. It is denoted
h= f g
or less rigorously h(x) = f (x) g(x). Note that the convolution of two
functions does not always exist.
Example 7.65 The convolution of the Heaviside function with itself is given by
H H (x) =
H (t) H (x t) dt = x H (x).
Conversely, the convolution of H with x 7 1/
|x| is not defined.
Exercise 7.4 Show that the convolution product is commutative, that is, h = f g = g f
whenever one of the two convolutions is defined.
Exercise 7.5 Let a, b R+ be real numbers such that a 6= b . Compute the convolution of
x 7 e |ax| with x 7 e |b x| .
(Solutions page 245)
Distributions I
212
It is possible to interpret graphically the convolution of two functions.

Indeed, denoting by h = f g their tensor product, we have
Z
Z
f g(x) = f (t) g(x t) dt = h(t, x t) dt,
which means that one integrates the function h : R2 C on the path formed
by the line D x with slope 1 passing through the point (0, x).
Example 7.66 Let f = and g(x) = (x 1 ), where is the rectangle function (see
2
page 213). Put h = f g and k = f g. Then the value of k(x) at a point x R is

given by the integral of h on the line D x below. The gray area corresponds to those points
where h(x, y) = 1, and the white area to those where h(x, y) = 0.
(0, x)
Dx
x
(x, 0)
If the diagonal line D x intersects the support of f g in a finite segment,

then the convolution of f by g is well defined at the point x.
Thanks to this interpretation, we have therefore proved the following result:
THEOREM 7.67 The convolution ot two locally integrable functions f and g exists
whenever one at least of the following conditions holds:
i) the functions f and g both have bounded support;

ii) the functions f and g are both zero for sufficiently large x (they have bounded
support on the right);
iii) the functions f and g are both zero for sufficiently small (negative) x (they have
bounded support on the left).
Those conditions are of course sufficient, but not necessary.
(where, as defined
Exercise 7.6 Show graphically that H H is well defined, but not H H
(x) = H (x)).
previously, H
Convolution
213
7.6.d Fuzzy measurement

The rectangle function is defined by
(x)
(x) =
1 if |x| <
0 if |x| >
1
2
1
,
2
x
21
1
2
Let now a > 0 be fixed. Denoting by fa the convolution of f by a similar

rectangle with width a and height 1/a, we have
Z
1 x 1 x+a/2
fa (x) = f (x)
=
f (t) dt,
a
a
a xa/2
which is simply the average of the function f on an interval around x of

length a. It is then possible to model simply an imperfect measurement of the
function f , for which the uncertainty is of size a. Then details of size a
disappear, whereas those of size a remain clearly visible.
This model can be refined by convolution of f by a gaussian or a lorentzian
(see Example 279), for instance.
Example 7.68 A spectral ray has the shape of a lorentzian (a classical result of atomic physics),
but what is observed experimentally is slightly different. Indeed, the atoms of the emitting gas
are not at rest, but their velocities are distributed according to a certain distribution, given by
a maxwellian function (this is the name given to the gaussian by physicists, because of the
Maxwell-Boltzmann distribution). Due to the Doppler effect, the wavelengths emitted by the
atoms are shifted by an amount proportional to the speed; therefore the spectral distribution
observed is the convolution of the original shape (the lorentzian) by a gaussian, with variance
proportional to the temperature.

It will be seen12 that the sequence of functions n(nx) nN converges (in
a sense that will be made precise in Definition 8.12 on page 230) to the Dirac
distribution . Thus, it can be expected that, in the limit [n ], since
the convolution does not change anything (the precision of measurement is
infinite), we have f = f ; thus we expect that will be a unit element for
the convolution product.
To show this precisely, we first have to extend the convolution to the
setting of distributions.
12
See Exercise 8.12 on page 242.
Distributions I
214
7.6.e
Convolution of distributions
By the same method as before, we wish to extend the convolution product

to distributions, so that it coincides with the convolution of functions in
the case of regular distributions. Now, for two locally integrable functions f
and g (such that the convolution f g is defined), we have for any D(R):
Z
f g, = ( f g)(t) (t) dt
Z
Z

= (t)
f (s) g(t s) ds dt
ZZ
=
f (x) g( y) (x + y) dx dy
(x = s, y = t s)

= f (x) g( y), (x + y) ,
since the jacobian in this change of variable is equal to | J | = 1.
DEFINITION 7.69 (Convolution of distributions) Let S and T be two distri-
butions. The convolution of S and T is defined by its action on any test

function D given by

def
S T , = S (x) T ( y), (x + y) .
The convolution of distributions does not always exist! The general conditions for its existence are difficult to write down. However, we may note that,
as in the case of functions, the following result holds:
THEOREM 7.70 The convolution of distributions with support bounded on the left
(resp. bounded on the right) always exists.

The convolution of an arbitrary distribution with a distribution with bounded
support always exists.
Indeed, if x 7 (x) is a nonzero test function, the map (x, y) 7 (x + y)

does not have bounded support, that is, it is not a test function on R2 ! If the
support of x 7 (x) is a segment [a, b], then the support of (x, y) 7 (x+ y)
is a strip as represented below:
b
a
Convolution
215
In particular, if we denote by the support of S and by the support

of T , then the convolution of S and T has a sense if the intersection of
an arbitrary strip of this type with the cartesian product is always
bounded. The reader can easily check that this condition holds in all the
situations of the previous theorem.
PROPOSITION 7.71 Let S and T be two distributions. Assume that the convolu-
tion S T exists. Then T S exists also and S T = T S (commutativity of

the convolution).
Similarly, if S , T and U are three distributions, then
S (T U ) = (S T ) U = S T U
if S T , T U , and S U make sense (associativity of the convolution product).

Counterexample 7.72 Let 1 denote as usual the constant function 1 : x 7 1. This function is
locally integrable and therefore defines a regular distribution, Moreover, the reader can check
that (1 ) = 0, hence (1 ) H = 0. But, on the other hand, ( H ) = , and hence
1 ( H ) = 1, which shows that
(1 ) H 6= 1 ( H ).
This can be explained by the fact that 1 H does not exist.
7.6.f
Applications
Let f and g be two integrable functions. They define regular distributions,

and, moreover (as easily checked),
f 1 and Rg 1 exist and are in fact conR
stant and equal to 1 f = f (x) dx, 1 g = g(x) dx, respectively. Assume
moreover that f g exists. Then, writing 1 [ f g] = [1 f ] g, we deduce
the following result:
THEOREM 7.73 Let f and g be two integrable functions such that the convolution
of f and g exists. Then we have

Z
Z
Z

[ f g](x) dx =
f (x) dx
g(x) dx .
Remark 7.74 The convolution product is the continuous equivalent of the Cauchy
product
P
P of
absolutelyP
convergent series. Recall that the Cauchy product of the series
the series wn such that
n
P
wn =
ak b nk .
an and
b n is
k=0
By analogy, one may write w = a b . Then we have
X
X

X
(a b )n =
an
bn .
n=0
n=0
n=0
Using the definition of convolution, it is easy to show the following relations:
Distributions I
216
THEOREM 7.75 Let T be a distribution. Then we have
T = T = T.
For any a R, the translate of T by a can be expressed as the convolution

(x a) T (x) = T (x) (x a) = T (x a),
and the derivatives of T can be expressed as

T = T = T
and (m) T = T (m) = T (m)
for any m N.
Proof
Let T D . Then, for any D, we have

T , = T (x) ( y), (x + y)
D

= T (x), ( y), (x + y) = T (x), (x) = T , x ,
which shows that T = T . Hence we also have T = T .
Let a R. Then
D

T a , = T (x), ( y a), (x + y) = T (x), (x + a) = T (x a), (x) ,
which shows that T a = Ta .
Finally, for any D,

D

T , = T (x), ( y), (x + y) = T (x), (x) = T , .
It follows that T = T = T .
In addition, if T = R S , then T = T = R S = R S .
Similarly, T = T = R S , which proves the following theorem:
THEOREM 7.76 To compute the derivative of a convolution, it suffices to take the
derivative one of the two factors and take the convolution with the other; in other
words, for any R, S D , we have
[R S ] = R S = R S .
when R S exists. In the same manner, in three dimensions, for any S , T D (R3 )
such that S T exists, we have
(S T ) = S T = S T .
7.6.g
The Poisson equation
We have admitted that the electrostatic potential V created by a distribution

of charge is given by

1
1
1
1
V ( r) =
( r ),
=
40
k r rk
40
k rk
Physical interpretation of convolution operators
217
(if this expression makes sense), where k r k denotes, with a slight abuse of
notation, the function r 7 k rk. It is then easy to compute, in the sense of
distributions, the laplacian of the potential. Indeed, it suffices to compute the
laplacian of one of the two factors in this convolution. Since we know that
(1/r) = 4, we obtain

1
1
1
1
1
V =

(4)
=
=
40
k rk
40
k rk
40
= .
0
Thus we have recovered a classical result:

THEOREM 7.77 (Poisson equation) The laplacian of the electrostatic potential cre-
ated by a distribution of charge is
V = .
0
7.7
We are going to see now that many physical systems, in particular measuring equipment can be represented by convolution operators. Suppose we
have given a physical system which, when excited by an input signal E depending on time, produces an output signal, denoted S = O (E ). We make
the assumptions that the operator O that transforms E into S is
linear;
continuous;
invariant under translations, that is, the following diagram commutes

(this means that composing the arrows from the upper left corner to
the bottom right corner in the two possible ways has the same result):
E (t)
translationy
S (t)
ytranslation
E (t a) S (t a)
Distributions I
218
THEOREM 7.78 If the three conditions stated above hold, then the operator O is a
convolution operator, that is, there exists R D such that
for any E .
S = O (E ) = E R = R E
The converse is of course true.

Remark 7.79 The case of an excitation depending on time is very simple. If, on the other
hand, we consider an optical system where the source is an exterior object and where the
output is measured on a photo-detector, the variable corresponding to the measurement (the
coordinates, in centimeters, on the photographic plate) and the variable corresponding to the
source (for instance, the angular coordinates of a celestial object) are diffent. One must then
perform a change of variable to be able to write a convolution relation between the input and
output signals.
Knowing the distribution R makes it possible to predict the response of

the system to an arbitary excitation. The distribution R is called the impulse
response because it corresponds to the output to an elementary input:
O
O () = R = R.
Note that often, in physics, the operator linking S (t) and E (t) is a differential operator.13 One can then write D S = E where D is a distribution
which is a combination of derivatives of the unit .14 The distribution R,
being the output corresponding to a Dirac peak, therefore satisfies
R D (t) = D R (t) = (t).
|{z}
E (t)
Since is the neutral element for the convolution, the distribution R is

thus the inverse of D for the convolution product. This point of view is
developed in Section 8.4 on page 238. In general, this operator D is given
by physical considerations; we will describe below the technique of Green
functions (Chapter 15), where the purely technical difficulty is to compute, if
it exists, the convolution inverse of D .
Remark 7.80 Experimentally, it is not always easy to send a Dirac function as input. It is
doable in optics (use a star), but rather delicate in electricity. It is sometimes simpler to send
a signal close to a Heaviside function. The response to such an excitation, which is called the
step response, is then S = H R. If we differentiate this, we obtain S = H R = R = R,
which shows the following result:
THEOREM 7.81 The impulse response is the derivative of the step response.
13
Consider, for instance, a particle with mass m and position x(t) on the real axis, subject
to a force depending only of time F (t). Take, for example, the excitation E (t) = m1 F (t) and
the response S (t) = x(t). The evolution equation of the system is S (t) = E (t).
14
In the previous example, one can take D = , since S = S .
219
through the telescope
A star
Two stars
A window
Fig. 7.3 Some objects to look at with a telescope.

If T is sufficiently smooth, this differentiation can be performed numerically; however, differentiation is often hard to implement numerically. This is the price to pay.
DEFINITION 7.82 (Causal system) A system represented by a temporal con-
volution operator (i.e., by a time variable) is called a causal system if its

impulse response R(t) is zero on R (which means that R, = 0 for any
test function which is zero for t 0).
In a causal system, the effect of a signal cannot come before its cause
(which is usually expected in physics).
Example 7.83 (in optics) Look at the sky with a telescope. Each star is a point-like light
source, modeled by a Dirac . Represent the light intensity of a star centered in the telescope
by a distribution a; then the image of this star after passing through the telescope is given
by a function of the intensity a(x, y), showing the diffraction spot due to the finite size of
the mirror and to the Newton rings.15 Any incoherent luminous object (which is indispensable
to ensure linearity with respect to the light intensity functions), with intensity function in the
coordinates (x, y) equal to F (x, y), will give an image F (x, y) (x, y) (see Figure 7.3).
15
The function is studied in Chapter 11, page 320.
Distributions I
220
7.8
Discrete convolution
It should be noted that it is quite possible to define a discrete convolution
to model the fuzziness of a digitalized picture. Such a picture is given no
longer by a continuous function, but by a sequence of values n 7 xn with n
[[1, N ]] (in one dimension) or (m, n) 7 xmn with (m, n) [[1, N ]] [[1, M ]]
(in two dimensions). Each of these discrete numerical values defines a pixel.
In the one-dimensional case, we are given values n for n [[1, N ]] (and
we put n = 0 outside of this interval), and the convolution is defined by
y =x
with yn =
+
X
k xnk ,
k=
the support of y being then equal to [[1 M , N + M ]]. In two dimensions,

the corresponding formula is
y =x
with ymn =
+
X
+
X
k xmk,n .
k= =
Photo-processing softwares (such as Photoshop, or the GIMP, which is Free

as in Free Speech software) provide such convolutions. An example of the
result is given in Figure 7.4 on the facing page.
The discrete convolution may be written as a matrix operation, with the
advantages and inconveniences of such a representation. The inversion of this
operation is done with the inverse matrix when it exists. We refer the reader
desirous to learn more to a book about signal theory.
The exercises for this chapter are found at the end of Chapter 8, page 241.
Discrete convolution
Fig. 7.4 A square in Prague, before and after discrete convolution.
221
Laurent Schwartz, born in Paris in 1915, son of the first Jewish surgeon in the
Paris hospitals, grandson of a rabbi, was a student at the cole Normale Suprieure.
During World War II, he went to Toulouse and then to Clermont-Ferrand, before
fleeing to Grenoble in 1943. In 1944, he barely escaped a German raid. After the war,
he taught at the University of Nancy, where the main part of the Bourbaki group
was located. He came back to Paris in 1953 and obtained a position at the cole
Polytechnique in 1959.
He worked in particular in functional analysis (under the influcence of
Dieudonn) and, in 1944, created the theory of distributions. Because of these
tremendously important works, he was awarded the Fields Medal in 1950. He showed
how to use this theory in a physical context.
Finally, one cannot omit, side by side with the mathematician, the political militant, the pacifist, who raised his voice against the crimes of the French state during
the Algerian war and fought for the independence of Algeria [82].
Laurent Schwartz passed away on July 4, 2002.
Chapter
Distributions II
In this chapter, we will first discuss in detail a particular distribution which is very
useful in physics: the Cauchy principal value distribution. Notably, we will derive
the famous formula
1
1
= pv i,
x i
x
which appears in optics, statistical mechanics, and quantum mechanics, as well as
in field theory. We will also treat the topology on the space of distributions and
introduce the notion of convolution algebra, which will lead us to the notion of
Green function. Finally, we will show how to solve in one stroke a differential
equation with the consideration of initial conditions for the solution.
8.1
Cauchy principal value
8.1.a
1
x
Definition
The function x 7 does not define a regular distribution, since it is not

integrable in the neighborhood of x = 0. On the other hand, one can define,
for D(R), the limit
Z +
Z
(x)
(x)
def
pv
dx = lim+
dx
0
x
|x| x
224
Distributions II
C + ()
x0
Fig. 8.1 The contour + .
which provides the definition of the distribution pv 1x , given by

Z +
Z
Z

1
(x)
(x)
(x)
def
pv , = pv
dx = lim+
dx +
dx .
0
x
x
x
This is then generalized as follows:
DEFINITION 8.1 Let x0 R. The Cauchy principal value pv
1
xx0
is defined
as the distribution on R given by

Z
Z
1
(x)
(x)
def
pv
, = pv
dx = lim+
dx.
0
x x0
x x0
|xx0 |> x x0
8.1.b Application to the computation of certain integrals

The situation that we are considering is the computation of an integral of the
type
Z +
f (x)
dx,
with f : R C.
x x0
This integral is not correctly defined without additional precision; indeed, if
f (x0 ) 6= 0, the function x 7 f (x)/(x x0 ) is not Lebesgue integrable. One
possibility is to define the integral in the sense of the principal value.
We will assume that f is the restriction to the real axis of a function
meromorphic on C and holomorphic at x0 . Moreover, we will assume that f
decays sufficiently fast at infinity to permit the use of the method of residues.
For instance, we will require that this decay holds at least in the upper halfplane.
Now define a contour + as in Figure 8.1: denote by C + () the small halfcircle with radius centered at x0 and above the real axis. Letting the radius
of the large circle go to infinity (and assuming the corresponding integral goes
to 0), we have then
Z
Z
Z +
Z
f (z)
f (x)
f (x)
f (z)
dz =
dx +
dx +
dz ;
x x0
+ z x0
x x0
C + z x0
225
x0
C ()
Fig. 8.2 The contour .
but the last integral is equal to

Z0
f ( e i )
i e i d = i f (x0 ) + O()
i
[ 0+ ].
Thus, taking the limit [ 0+ ], we obtain

Z +
X
f (z)
f (x)
2i
Res
; p = pv
dx i f (x0 ),
(8.1)
z x0
x x0
pH
P
where the sum pH is the sum over the poles of f located in the open
def
upper half-plane H = {z C ; Im(z) > 0}.

So, to deal with the integral on a contour on which a pole is located,
it was necessary to deform this contour slightly in this case, by making a
small detour around the pole to avoid it. One may wonder what would have
happened, had another detour been chosen. So take the contour defined
as in Figure 8.2. We then have
Z
Z +
Z
f (z)
f (x)
f (z)
dz = pv
dx +
dz
z x0
x x0
C z x0
Z +
f (x)
= pv
dx + i f (x0 ).
x x0
Since the integral on the left-hand side is equal to 2i times the sum of the
residues located in the upper half-plane, plus the residue at x0 , which is equal
to 2i f (x0 ), the final formula is the same as (8.1).
8.1.c Feynmans notation

When evaluating an integral on a contour that passes through a pole, it must
be specified on which side the pole the pole will be avoided (the result depends
only on the side that is chosen, not on the precise shape of the detour).
Rather than going around the pole, Richard Feynman proposed modifying the initial problem by moving the pole slightly off the real axis while
Distributions II
226
keeping the original contour intact. Suppose the function considered has a
pole on the real axis. Then one can write the symbolic equivalences
and
which means that the pole has been moved in the first case to x0 i and in
the second case to x0 + i, and then the limit [ 0+ ] is taken at the end of
the computations.
To see clearly what this means, consider the equalities
Z
Z
Z
f (z)
f (z)
f (z)
dz = lim+
dz,
dz = lim+
0
0
+ z x0
+ z x0 + i
z x0 + i
where is the undeformed contour. Indeed, the first equality is a consequence
of continuity under the integral sign, and the second equality comes from
Cauchys theorem, which, for any given > 0, justifies deforming + into ,
since no pole obstructs this operation. Thus we get, performing implicitly the
operation lim+ ,
Z
f (z)
dz =
z x0
f (x)
dx = pv
x x0 + i
f (x)
dx i f (x0 ),
x x0
or, in a rather more compact fashion

1
1
, f = pv
, f i (x x0 ), f .
x x0 + i
x x0
Symbolically, we write:
1
= pv
x x0 + i
1
x x0
i (x x0 ),
and this formula must be applied to a meromorphic function f which is

holomorphic on the real axis.
Similarly, we have the following result:

1
1
= pv
+ i (x x0 ).
x x0 i
x x0
227
Notice and this will be convenient in certain computations that

1
1
pv
= Re
.
x x0
x x0 i
8.1.d Kramers-Kronig relations

Let F be a function from C into C which is meromorphic on C but holomorphic on the closed upper half-plane, which means that all the poles of F have
strictly negative imaginary part. Consider now the contour C described below,
with an arbitrarily large radius:
From Cauchys formula, we have
Z
F ( )
1
F (z)
d =
0
2i C z
if Im(z) > 0,
if Im(z) < 0.
In this formula, put z = x + i with x R and > 0:

Z +
1
F (x )
F (x + i) =
dx .
2i x x i
Assume that the integral on the half-circle tends to 0 when the radius tends to
infinity. By letting tend to 0, and by continuity of F , the preceding results
lead to
Z +
F (x ) 1
1
pv
dx + F (x),
F (x) =
2i
2
x x
hence
Z +
F (x )
1
F (x) =
pv
dx .
i
x x
This equation concerning F becomes much more interesting by taking successively its real and imaginary parts; the following theorem then follows:
THEOREM 8.2 (Kramers-Kronig relations) Let F : C C be a meromorphic
function on C, holomorphic in the upper half-plane and going to 0 sufficiently fast at
infinity in this upper half-plane. Then we have

Z +

Im F (x )
1
Re F (x) = pv
dx
x
228
and
Distributions II

1
Im F (x) = pv

Re F (x )
dx
x x
These are called dispersion relations and are very useful in optics (see, e.g., the
book by Born and Wolf [14, Chapter 10]) and in statistical physics. Physicists
call these formulas the Kramers-Kronig relations,1 while mathematicians say
that Re(F ) is the Hilbert transform of Im(F ).
Remark 8.3 If F is meromorphic on C and holomorphic in the lower half-plane (all its poles
have positive imaginary parts), then it satisfies relations which are dual to those just proved
(also called Kramers-Kronig relations, which does not simplify matters):

Z +
Z +

Im F (x )
Re F (x )
1
1
dx and Im F (x) = pv
dx .
Re F (x) = pv
x x
x x
We will see, in Chapter 13, that the Fourier transform of causal functions t 7 f (t) (those
that vanish for negative values of the variable t), when it exists, satisfies the Kramers-Kronig
relations of Theorem 8.2.
Remark 8.4 What happens if F is holomorphic on both the upper and the lower half-plane?
Since it is assumed that F has no pole on the real axis, it is then an entire function. The
assumption that the integral on a circle tends to zero as the radius gets large then leads (by the
mean value property) to the vanishing of the function F , which is the only way to reconcile
the previous formulas with those of Theorem 8.2.
Remark 8.5 In electromagnetism, the electric induction D and the electric field E are linked,
for a monochromatic wave, by
D( x, ) = () E( x, ),
where () is the dielectric constant of the material, depending on the pulsation of the
waves. This relation can be rewritten (via a Fourier tranform) in an integral relation between
D( x, t) and E( x, t):
Z +
D( x, t) = E( x, t) +
G () E( x, t ) d,
where G () is the Fourier transform of () 1. In fact, since the electric field is the physical
field and the field D is derived from it, the previous relation must be causal, that is, G () = 0
for any < 0. One of the main consequences is that the function 7 () is analytic in the
lower half-plane of the complex plane (see Chapter 12 on the Laplace transform). From the
Kramers-Kronig relations, interesting information concerning the function 7 () can be
deduced. Thus, if the plasma frequency is defined by p 2 = lim 2 [1 ()], we obtain
the sum rule

2
Im () d.
0
The reader is invited to read, for instance, the book by Jackson [50, Chapter 7.10] for a more
detailed presentation of this sum rule.
p 2 =
Dispersion relations first appeared in physics in the study of the dielectric constant of
materials by R. de L. Kronig in 1926 and, independently, in the theory of scattering of light
by atoms by H. A. Kramers in 1927.
229
Exercise 8.1 Using the Kramers-Kronig relations, show that
sin x
dx = .
x
2
(Solution page 245)
8.1.e A few equations in the sense of distributions

Recall (Theorem 7.28 on page 191) that the solutions of the equation xT (x) =
0 are the multiples of .
PROPOSITION 8.6 We have x pv
1
x
= 1 (where 1 is the constant function x 7 1).
Proof. Indeed, for any test function D(R), we have by definition of the product
of a distribution (pv 1x ) by a C function:
Z
Z +
1
1
x (x)
x pv , = pv , x = lim+
dx =
(x) dx = 1, .
0
x
x
x
|x|>
THEOREM 8.7 (Solutions of x T ( x) = 1) The solutions, in the space of distribu-
tions, of the equation
x T (x) = 1,
are the distributions given by T (x) = pv 1x + , with C.
Proof. If T satisfies x T (x) = 1, then S = T pv(1/x) satisfies x S (x) = 0 and
thus, according to Theorem 7.28 on page 191, S is a multiple of .
Theorem 7.28 generalizes as follows:

PROPOSITION 8.8 Let n N. The solutions of equation x n T (x) = 0 are the
linear combinations of , , . . . , n1 .
Let f be a function continuous on R, with isolated zeros. For each such

zero x i of f , we say that x i is of multiplicity m if there exists a continuous
function g, which does not vanish in a neighborhood of x i , such that f (x) =
x m g(x) for x R.
If f has only isolated zeros with finite multiplicity, then one can treat the
equation f (x) T (x) = 0 locally around each zero in the same manner as
x m T (x) = 0.
THEOREM 8.9 (Solutions of f ( x) T ( x) = 0) Let f be a continuous function
on R. Denote by Z( f ) = {x i i E } the set of its zeros, and assume these zeros
are isolated and with finite multiplicities m i .
The distributions solutions to the equation f (x) T (x) = 0 are given by
T (x) =
i 1
X mX
iI k=0
where i,k are arbitrary complex numbers.
i,k (k) (x x i ),
Distributions II
230
Example 8.10 The equation sin x T (x) = 0 has for solutions the distributions of the form
nZ
n (x n).
The equation (cos x 1) T (x) = 0 has for solutions the distributions of the form

P
n (x n) + n (x n) .
nZ
THEOREM 8.11 (Solutions of ( x 2 a 2 ) T ( x) = 1) Let a > 0 be given. The
distribution solutions of the equation (x 2 a 2 ) T (x) = 1 are those given by

pv
x2
1
+ (x a) + (x + a),
a2
where and are complex numbers, and where we denote
1
1
1
1
pv 2
=
pv
pv
.
x a2
2a
x +a
x a
8.2
Topology in D
8.2.a
Weak convergence in D
DEFINITION 8.12 A sequence of distributions (Tk )kN converges weakly (or

simply converges) in D if, for any test function D, the sequence of
complex numbers Tk , converges in C. We then denote by T , this limit,
which defines a linear function T on D, called the weak limit of (Tk )kN
in D .
THEOREM 8.13 If (Tk )kN converges weakly in D to a functional T , then T is
continuous and is therefore a distribution. The weak convergence in D is denoted by

D
Tk T .
Exercise 8.2 For n N, let Tn be the regular distribution associated to the locally integrable
function x 7 sin(nx). Show that the sequence (Tn )nN converges weakly to 0 in D .
THEOREM 8.14 (Continuity of differentiation)
in
D (R)
have
(m)
Tk
The operation of differentiation

D
is a continuous linear operation: if Tk T , then for any m N, we

D
T (m) .
Obviously this extends to partial derivatives of distributions in Rn .
Topology in D
231
Fig. 8.3 The sequence of gaussian functions (gn )nN , for n = 1, . . . , 5.
This theorem is extremely powerful and justifies the interchange of limits

and derivatives, which is of course impossible in the sense of functions.2
8.2.b Sequences of functions converging to

DEFINITION 8.15 A Dirac sequence (of functions) is any sequence of functions fk : R R locally integrable, such that
there exists a real number A > 0 such that, for any x R and any
k N, we have

|x| A = fk (x) 0 ,
for any a R, a > 0,

Z
fk (x) dx 1
|x|a
and
cv.u.
fk (x) 0
1
on a < |x| < .
a
Example 8.16 An example of a Dirac sequence is the sequence ( fn )nN defined by fn (x) =
n(nx) for any x R and any n N. (Check this on a drawing.)
Example 8.17 The sequence of functions (gn )nN defined by
n
def
2 2
gn (x) = p e n x
for x R
Consider, for instance, the sequence of functions ( fn )nN defined by

n
X
1
fn (t) =
cos 2k t
for any t R.
2k
k=1
The functions fn are all differentiable, but their limit f is a sawtooth function, which is not
differentiable at the points of the form n with n Z.
Distributions II
232
Fig. 8.4 The sequence of functions n (x) =
sin2 (nx)
for n = 1 to 5.
n2 x 2
is a Dirac sequence (see Figure 8.3 on the previous page).
THEOREM 8.18 Any Dirac sequence converges weakly to the Dirac distribution
in D (R).
Proof. See Appendix D.

Example 8.19 Similarly, the sequence of functions (n )nN defined by
sin2 (nx)
for x R
n2 x 2
is a Dirac sequence (see Figure 8.4) and thus converges also to .
n (x) =
Remark 8.20 Without the positivity assumption for x close to 0, there is nothing to prevent
contributions proportional to to arise in the limit! Thus, consider the sequence of functions
(kn )nN given by the graph
n2
1
n
n1
n2
D
It is easy to show that kn . Let ( fn )nN be a Dirac sequence; then the sequence
( fn + kn )nN , which still satisfies , but not , converges to + and not to .
Topology in D
233
10
8
6
4
2
0
0,5
0,5
Fig. 8.5 The functions b n : x 7
sin 2nx
, for n = 1, . . . , 5.
x
Consider now the sequence of functions

bn (x) =
sin 2nx
.
x
All the functions bn share a common envelope x 7 1/x and the sequence
of real numbers (bn (x))nN does not converge if x
/ Z (see Figure 8.5). We do
have bn (0) , but, for instance,

1
(1)k /2 if n = 2k + 1,
bn
=
0
if n = 2k.
4
In particular, this sequence (bn )nN is not a Dirac sequence. Yet, despite this
we have the following result:
PROPOSITION 8.21 The sequence (bn )nN converges to in the sense of distribu-
tions.
How to explain this phenomenon?

Proving the convergence of (bn )nN to is, because of the lack of positivity, harder than for a Dirac sequence. In the case that concerns us, the
sequence of functions (bn )nN exhibits ever faster oscillations when we let n
tend to infinity (see Figure 8.5), and, after getting rid of the singularity at 0,
we can hope to be able to invoke the Riemann-Lebesgue
lemma (see TheoR
rem 9.44). Moreover, it is easy to compute that bn = 1 for all n N (for
example, by the method of residues).
Proof of Proposition 8.21. Let D be a test function and let [M , M ] be an
interval containing the support of .
Distributions II
234
We can always write (x) = (0) + x (x), where if of C class. Then we have
ZM
ZM
sin 2nx
dx +
(x) sin(2nx) dx,
b n , = (0)
x
M
M
and the first integral, via the substitution y = nx, tends to the improper integral
Z +
sin x
dx = 1,
x
whereas the second integral tends to 0, by the Riemann-Lebesgue lemma.
Finally, the convergence to is highly useful because of the following

theorem:
THEOREM 8.22 Let ( fn )nN be a Dirac sequence of functions. If g is a bounded
continuous function on R, then fn g exists for all n N and the sequence ( fn g)nN
converges uniformly to g.
If g is not continuous but the limits on the
left or the right of g exist at the point x,
then fn g(x) tends to 12 g(x ) + g(x + ) .
8.2.c
Convergence in D and convergence

in the sense of functions
The previous examples show that a sequence ( fn )nN of regular distributions

may very well converge in D to a regular distribution f without converging
pointwise almost everywhere (and even more without converging uniformly)
in the sense of functions. In particular, the Riemann-Lebesgue lemma is
equivalent to the following proposition:
PROPOSITION 8.23 The sequence of regular distributions ( fn )nN defined by the
functions fn : x 7 sin(nx) converges to 0 in D , but it converges pointwise only for

x Z.
8.2.d Regularization of a distribution

What happens when we take the convolution of a function f of C 1 class
with a function g of C 2 class? Assuming that this convolution exists and
denoting h = f g, the function h will be infinitely differentiable in the sense
of distributions, but how many times will it be differentiable in the sense of
functions?
Compute the derivative. It is possible to write h = f g and, since f
and g are functions, h is well-defined in the sense of functions. Continue. It
is not possible to differentiate f more than that, but g is still available, and
so h = f g is still a function, as is the third derivative h = f g . But
going farther is not possible.
Similarly, the convolution of a function of C k class with a function of C
class will produce a function of C k+ class.
Topology in D
235
These results generalize to distributions. Of course, it is not reasonable

to speak of C class (since distributions are always infinitely differentiable
anyway). But regularity for a distribution consists precisely in being what
is called a regular distribution, namely, a distribution associated to a locally
integrable function.
An interesting result (which is, in fact, quite intuitive) is that the convolution by a C function smoothes the object which is subject to the
convolution: the result is a function which is C . This is the case also if
we take the convolution of a distribution with a C function: the result is
a regular distribution associated to a function, and in certain cases one can
show that the result is also C .
THEOREM 8.24 (Regularization of a distribution) Let T D be a distribution
and C (R), and assume that T exists. Then T is a regular distribution

associated to a function f , which is given by

f (x) = T (x) = T (t), (x t) .
Let T D and D (so is now not only of C class, but also with
bounded support). Then T exists and is a regular distribution associated to a
function g of C class, with derivatives given by the formula

g (k) (x) = T (t), (k) (x t) .
8.2.e Continuity of convolution
The following theorem shows that the operation of convolution with a fixed
distribution is continuous for the weak topology on D :
THEOREM 8.25 (Continuity of convolution) Let T D . If (Sk )kN converges
weakly in D to T , and if moreover all the Sk are supported in a fixed closed subset
such that all the Sk T exist, then the sequence (Sk T ) converges in D to S T .
Consider now a Dirac sequence (n )nN , with elements belonging to D.
It converges weakly to in D (this is Theorem 8.18); moreover, we know
(Theorem 8.24) that if T is a distribution, we have n T C (R). Since
n T T , we deduce:
n
THEOREM 8.26 (Density of D in D ) Any distribution T D is the limit in

D of a sequence of functions in D. One says that D is dense in D .
In other words, any distribution is the weak limit of a sequence of functions of C class with bounded support.
Exercise 8.3 Find a sequence (n )nN in D which converges weakly to X in D .
(See Exercise 8.11 on page 242)
Distributions II
236
8.3
Convolution algebras
DEFINITION 8.27 A convolution algebra is any vector space of distributions
containing the Dirac distribution and for which the convolution product
of an arbitrary number of factors is defined.
There are two well-known and very important examples.
THEOREM 8.28 (E
E , D + and D ) The space E of distributions with bounded sup-
port is a convolution algebra.

The space of distributions with support bounded on the left, denoted D + and
the space of distributions with support bounded on the right, denoted D , are also
convolution algebras.
These algebras, as the name indicates, provide means of solving certain

equations by means of algebraic manipulations. In particular, many physical
problems can be put in the form of a convolution equation of the type
AX =B
where A and B are distributions (B is often called the source) and the
unknown X is a distribution.3 The problem is to find a solution and to
determine if it is unique. So we are, quite naturally, led to the search for a
distribution denoted A 1 , which is called the convolution inverse of A or,
sometimes, Green function of A, which will be such that
A A 1 = A 1 A = .
Remark 8.29 This inverse, if it exists, is unique in the algebra under consideration. Indeed, suppose
there exist, in a same algebra, two distributions A and A such that
Then
and therefore
A A = A A =
and
A A = A A = .
A = A = A (A A ) = (A A) A = A = A
A = A .
For instance, notice that any linear differential equation with constant coefficients

d
dn
a0 + a1 + + an n X (t) = B(t)
dt
dt
can be reexpressed in the form of a convolution product

a0 + a1 + + an (n) X (t) = B(t).
Convolution algebras
237
The convolution equation is then solved by putting

X = A 1 B,
and this solution is unique.
To summarize, we can state the following theorem:
THEOREM 8.30 Let A be a convolution algebra. The convolution equation
AX =B
admits a solution in A for arbitrary B A if and only if A has a convolution

inverse A 1 in this algebra. In this case, the solution is unique and is given by
X = A 1 B.
Remark 8.31 It is important to notice that we have always spoken of the inverse of a distribu-
tion in a given algebra A . In fact, a distribution A may very well have a certain convolution
inverse in one algebra and another inverse in a different algebra. Take, for instance, the case of
a harmonic oscillator characterized by a convolution equation of the type
( + 2 ) X (t) = B(t),
where B(t) is a distribution characterizing the input signal. The reader is invited to check that
( + 2 )
1
H (t) sin t = .
Since the distribution 1 H (t) sin t is an element of D+ , it follows that, for any B D+ , that
is, for any input which started at a finite instant t R and was zero previously, the equation
has a solution in D+ which is unique and is given by
X (t) = B(t)
1
H (t) sin t
causal solution.
But it is also possible to check that
1
( + 2 ) H (t) sin t = ,
and therefore that 1 H (t) sin t is a convolution inverse of ( + 2 ) in D . If the

input signal B(t) belongs to D , it is therefore necessary, to get the solution, to carry out the
convolution of B by 1 H (t) sin t:
X (t) = B(t)
1
H (t) sin t
anti-causal solution.
We also deduce that this convolution operator does not have an inverse in E .
One of the great computational problems that can confront a physicist

is to find the convolution inverse of an operator (its Green function). For
instance, the Green function of the Laplace operator in R3 is given by the
Coulomb potential 1/r, as we have seen by performing the relevant convolution product. But how, given an operator, can we find its inverse? If the
Distributions II
238
convolution operator is a differential operator (more precisely, a differential

operator applied to ), there exist tools which are particularly adapted to computing its inverse, namely the Fourier and Laplace transforms (see Chapters 10
and 12). Infinitely many technical details for the inversion of the laplacian
(with refinements thereof) are given in the reference [24], which is somewhat
hard to read, however.
8.4
Solving a differential equation
with initial conditions
Most students in physics remember very well that the evolution of a system
is given by a differential equation, but forget that the real problem is to solve
this equation for given initial conditions. It is true that, in the linear case,
these initial conditions are often used to parameterize families of solutions;
however, this is not a general rule.4
The Cauchy problem (an equation or system of differential equations,
with initial conditions) is often very difficult to solve, as every one of you
knows. Sometimes, it is very hard to show that a solution exists (without even
speaking of computing it!), or to show that it is unique.5
In the cases below, we have autonomous linear differential equations with
constant coefficients. The theoretical issues are therefore much simpler, and we
are just presenting a convenient computational method. Some cases which are
more complicated (and more useful in physics) will be presented on page 348.
8.4.a
First order equations
Consider a differential equation of the type

u + u = 0,
(e)
where we only wish to know u(t) for positive values of t, and u : R+ R is

a function satisfying the initial condition
u(0) = u0 R.
It is possible to use the tools of the theory of distributions to solve this
equation with its initial conditions.
4
The most spectacular counterexamples are given by chaotic systems.

For instance, the old problem of an electric charge in an electromagnetic field is only
partly solved [11]. The difficulty of this system is that it is described by highly nonlinear
equations where the charge influences the fields, which in turn influence the charge. The
reader interested in these problems of mathematical physics can read [84].
5
Solving a differential equation with initial conditions
239
We then slightly change our point of view, by looking no longer for a

solution in the sense of functions to this differential equation, but for a
solution in the sense of distributions, of the type U (t) = H (t) u(t). (This
solution will then be zero for negative time, which is of no importance since
we are interested in positive time only.)
If u is a solution of (e), which equation will the distribution U satisfy?
Since U = u + u0 , we see that U satisfies the equation U + U = u0 ,
which may be put in the form of a convolution
[ + ] U = u0 .
(E)
There only remains to find the unique solution to this equation (in the sense
of distributions) which lies in the convolution algebra D+ of distributions
with support bounded on the left.
is
THEOREM 8.32 The convolution inverse of [ + ] in D+
G (t) = H (t) e t .
(The inverse in D is t 7 H (t) e t .)
Proof. Indeed, we have G (t) = e t +(t), which shows that [ +]G = .
Taking the convolution of (E) with G (t) = H (t) e t , we find that

U = G (u0 ),
or
U = u0 H (t) e t .
It may seem that we used something like a hammer to crush a very small
fly. Quite. But this systematic procedure becomes more interesting when
considering a differential equation of higher order, or more complicated initial
conditions, for instance on surfaces (see Section 15.4, page 422).
8.4.b
The case of the harmonic oscillator
Consider now the case of a harmonic oscillator, controlled by the equation

d2 u
+ 2 u = 0,
dt 2
(e )
where we seek a function u(t) satisfying this differential equation for t 0,

with given initial conditions u(0) and u (0). For this purpose, we will look
for a distribution U (t), of the type U (t) = H (t) u(t), which satisfies the
differential equation in the sense of functions. What differential equation in
the sense of distributions will it satisfy? This will be, taking the discontinuities
into account,

+ 2 U = u(0) + u (0) .
(E )
Distributions II
240
THEOREM 8.33 (Harmonic oscillator) The convolution inverse of [ + 2 ]
in D+ is
1
H (t) sin t.
The convolution inverse of [ + 2 ] in D (see Remark 8.31 on page 237) is

G + (t) =
G (t) =
1
H (t) sin t.
Remark 8.34 There exist methods to rediscover this result, the proof of which is immediate:
such convolution inverses, or Green functions, of course do not come out of thin air. These
methods are explained in Section 15.2 on page 409.
Then we find directly, by taking the convolution of equation (E ) with G + ,

that

1
U (t) = H (t) sin t u(0) (t) + u (0) (t)

u (0)
= H (t) u(0) cos t +
sin t .
To summarize, the same problem can be seen from two different angles:
1. the first, classical, shows a differential equation with given initial conditions at the point 0;
2. the second shows a differential equations valid in the sense of distributions, without conditions at 0, for which a solution is sought in a given
convolution algebra (for instance, that of distributions with support
bounded on the left).
The right-hand side of (E ) shows what distribution is necessary to move

the system at rest (for t < 0) instantaneously to the initial state u(0), u (0) .
8.4.c Other equations of physical origin
We will see in the following chapters, devoted to the Fourier and Laplace
transforms, how to compute the convolution inverses (also called the Green
functions) of a differential operator.
Here we merely state an interesting result:
THEOREM 8.35 (Heat equation) The convolution inverse of the heat operator
in D+ , called the heat kernel, is
2
c
2
t
x
H (t) c x 2 /4t
G (x, t) = p
e
,
c t
Exercises
that is,
or
241

c (t) (x) G (x, t) = (2) (x, t),
c
2G
G
2 = (x) (t).
t
x
Proof. This result is proved ab initio in Section 15.4, page 422.
EXERCISES
Examples of distributions
Exercise 8.4 (Cauchy principal value) Show that the map defined on D by
1
pv : D 7 lim
0
x
|x|>
(x)
dx
x
which defines the Cauchy principal value is the derivative, in the sense of distributions, of
the regular distribution associated to the locally integrable function x 7 log |x|.
Exercise 8.5 (Beyond principal value) Show that the map
7 lim+
0
Z
(0)
(x)
dx
+ (0) log
2
x
defines a distribution, that is, that it is linear and continuous.
Exercise 8.6 (Finite part of 1/ x 2 ) We denote by fp(1/x 2 ) the opposite of the distribution
derivative of the distribution pv(1/x) (fp for finite part):
1 def
1
.
fp 2 = pv
x
x
Show that
Z +
Z
1
(x) (0)
fp 2 , =
log |x| (x) dx = lim+
dx,
0
x
x2
|x|>
and that, if 1 denotes as usual the constant function equal to 1, or the associated distribution,
we have
1
x 2 fp 2 = 1 .
x
Generalize these results to define distributions related to the (not locally integrable) functions
x 7 1/x n for any n N.
Exercise 8.7 Show that we have, for all a R,
(a x) =
whereas, if a > 0,
1
(x),
|a|
H (a x) = H (x).
Distributions II
242
Exercise 8.8 Prove the relation : x = . Compute similarly x . Generalize by com-
puting f (x) , where f is a (locally integrable) function of C class. Is it possible to relax

those assumptions?
Composition of a distribution and a function

Exercise 8.9 (Special relativity) We consider the setting of special relativity, with the momentum vector p = (E , p) = (p 0 , p), with a Minkowski metric of signature (+, , , ) and
the convention c = 1. Thus, we have p2 = E 2 p 2 .
Recall that for a classical particle, we have p2 = m 2 , where m is the mass of the particle. In
field theory, one is often led to compute integrals of the type
Z
(p2 m 2 ) H (p 0 ) f (p) d4 p.
What is the interpretation of the functions and H ? These integrals have the nice property
of being invariant under Lorentz transformations. Show that, if we put
p
def
E ( p) = p 2 + m 2 ,
then we have
(p2 m 2 ) H (p 0 ) d4 p =
d3 p
,
2E ( p)
which is the usual way of writing that for any sufficiently smooth function (indicate precisely
the necessary regularity) we have
Z
Z
d3 p
(p2 m 2 ) H (p 0 ) f (p) d4 p = f E ( p), p
.
2E ( p)
The important point to remember is the following:
THEOREM 8.36 The differential element d3 p/2E , where E 2 = p 2 + m 2 , is Lorentz invariant.

Exercise 8.10 Let S denote the sphere centered at 0 with radius R. Can we write S in
p
the form (r R), where r = x 2 + y 2 + z 2 ? Can we write it (r 2 R 2 )?
Convergence in D
Exercise 8.11 Let n N and R. We consider the function
fn : R \ Z R,
1 sin2 nx/2
.
2n sin2 x/2
R +1
Show that fn can be extended by continuity to R. Compute 1 fn (t) dt and show that the
result is independent of n.
Show that the sequence ( fn )nN tends to X in the sense of distributions.
x 7
x 7
1 x D
(x)
a
a
[a 0+ ].
Exercises
243
Exercise 8.13 Let be a real number. Show that, in the sense of distributions, we have
X
nZ
sin |n| =
sin
.
1 cos
Differentiation in higher dimensions

Exercise 8.14 We consider a two-dimensional space-time setting. Define
E (x, t) =
1
0
if c t > |x| ,
if c t < |x| ,
or, in short-hand notation, E (x, t) = H (c t x). Compute E in the sense of distributions,

where the dAlembertian operator is defined by
def
2
1 2
2.
2
2
c t
x
Exercise 8.15 In the plane R2 , we define the function E (x, t) = H (x) H (t), and its asso-
ciated regular distribution, also denoted E . Compute, in the sense of distributions,

and E .
2
E
x t
Exercise 8.16 (The wages of fear) We model a truck and its suspension by a spring with
strength k carrying a mass m at one of its extremities. The wheels at the bottom part of the
spring are forced to move on an irregular road, the height of which with respect to a fixed
origin is described by a function of the distance E (x). The horizontal speed of the truck is
assumed to be constant equal to v, so that at any time t, the wheels are at height E (v t). Denote
by R(t) the height of the truck, compared to its height at rest (it is the response of the system).
i) Show that

d2
k
2
+
R(t) = 2 E (v t),
where 2 = .
dt 2
m
ii) In order to generalize the problem, we assume that E is a distribution. What is then
the equation linking the response R(t) of the truck to a given road profile E (x) D+ ,
where D+ is the space of distributions with support bounded on the left? Comment
on the choice of D+ .
iii) Compute the response of the truck to the profile of a road with the following irregularities:
a very narrow bump E (x) = (x);
a flat road E (x) = H (x);
a wavy road E (x) = H (x) sin(x/).
Compute the amplitude of the movements of the truck. How does it depend on the
speed?
iv) In the movie The Wages of Fear by Henri-Georges Clouzot (1953, original title Le salaire
de la peur), a truck full of nitroglycerin must cross a very rough segment of road. One
of the characters, who knows the place, explains that the truck can only cross safely at
very low speed, or on the contrary at high speed, but not at moderate speed. Explain.
Distributions II
244
PROBLEM
Problem 3 (Kirchhoff integrals) The wave equation in three dimensions, for a scalar quan-
tity , is
def
1 2
c2 t2
( r, t) = 0,
where we denote r = (x, y, z). Monochromatic waves with pulsation are given by
( r , t) = ( r) e i t .
1. Show that the complex amplitude then satisfies the equation called the Helmoltz
equation:

def
+ k 2 ( r) = 0,
with k = .
c
Let S be an arbitrary closed surface. We are looking for a distribution ( r ) such that
( r) inside S ,
( r ) =
0
outside S ,
and which is a solution in the sense of distributions of the Helmoltz equation.
2. Show that differentiating in the sense of distribution leads to

h
i

d
d
+ k 2 ( r) ( r ) =
S ( r) +
S ( r ),
dn
dn
where n is the interior normal to S (the choices of orientation followed here are those
in the book Principles of Optics by Born and Wolf [14]).
3. Show that a Green function of the operator [ + k 2 ], that is, its convolution inverse,
is given by
1 e i kr
def
G ( r) =
,
where r = k r k .
4 r
4. Deduce that
ZZ

d e i kq
e i kq d
1
d2 x
with q = k r xk .
( r ) =
4
dn q
q
dn
S
5. Interpretation: the complex amplitude inside a closed surface S , with given limit conditions, is the same as that produced by a certain surface density of fictive sources6
located on S . This result is the basis of what is called the method of images in electromagnetism and it is, in the situation considered here, the founding principle of the
scalar theory of diffraction in optics.
A more complete version of Kirchhoffs Theorem, for an arbitrary dependence on time,
is presented in the book of Born and Wolf [14, p. 420].
Possibly very difficult to compute explicitly! Note that this computation requires knowing
both the amplitude and its normal derivative with respect to the surface. In some cases only
we may hope to exploit this exact formula, where symmetries yield additional information.
245
SOLUTIONS
Solution of exercise 7.2 on page 198
x E
1 D
1
(x),
=
(0)
|a|
a
a |a|
1
, .
=
a |a|
R
Solution of exercise 7.4 on page 211. Put h(x) = f (t) g(x t) dt and make the substitution y = x t, with jacobian |dy/dt| = 1, which yields
Z
Z
h(x) = f (x y) g( y) dy = g( y) f (x y) dy = [g f ](x).
(a x), (x) =
Solution of exercise 7.5 on page 211. A simple but boring computation shows that
e |ax| e |b x| =
2
[a e b|x| b e a|x| ].
a2 b 2
Solution of exercise 8.1 on page 229. Notice that sin x = Im(e i x ) and that z 7 e i z is
holomorphic on C and decays in the upper half-plane (at least in the sense of Jordans second
lemma). Using then the Kramers-Kronig formula from Theorem 8.2, we obtain
Z
Z
1 + sin x
sin x
dx =
dx
x
2 x
0
Z Z +
1
sin x
= lim+
+
dx
(continuous function)
2 0
x
Z +
1
Ime i x
= pv
dx = Re(e i 0 ) = .
2
x
2
2
Solution of exercise 8.4. Let D. We then have

log |x| , = log |x| , (x) =
= lim
0+
log |x| (x) dx

Z

log |x| (x) dx,
since x 7 log |x| is locally integrable (the integral is well defined around 0). Integrating each
of the two integrals by parts, we get
Z

(x)
log |x| , = lim+ log() () log() () +
dx
0
x
|x|>
Z +
(x)
= pv
dx,
x

which shows that log |x| = pv 1x .
Solution of exercise 8.5. First, we show that the stated limit does exist. Let D. Integrating by parts twice, we get
Z +
Z +
(x)
()
dx
=
(x) log x dx () log +

,
x2
Distributions II
246
and therefore
Z +
(0)
(x)
dx
+ (0) log =
2
x
Z +

() (0)
(x) log x dx + (0) () log +

.
The second term tends to 0, and the third to (0) as [ 0+ ]. The first term is integrable
when = 0. Hence the expression given has a limit, which we denote by T , .
The linearity is obvious, and hence we have to prove the continuity, and it is sufficient to
show the continuity at 0 (the zero function of D).
Let (n )nN be a sequence of functions in D, converging to 0 in D. All the functions n
have support included in a fixed interval [A , A], where we can assume for convenience and
without loss of generality that A > e. Then
ZA
T , n =
n (x) log x dx + n (0),
0
and hence

T , n k k (A log A A) + k k .
n
n
Choose > 0. There exists an integer N such that for n N we have kn k /(A log A
A) and kn k , so that |T , n | 2 for any n N . This proves that T , n tends to 0
as [n ].
Since we have proved the continuity of T , this linear functional is indeed a distribution.
Solution of exercise 8.6. To generalize, define:

fp
(1)n1 dn
(1)n1 dn1
1
1
=
log |x| =
pv .
n
n
x
(n 1)! dx
(n 1)! dx n1
x
Solution of exercise 8.7. Let a 6= 0. Then for any D we have
1
1
(a x), =
(x), (x/a) =
(0) =
, .
|a|
|a|
|a|
For any a > 0, the functions x 7 H (x) and x 7 H (a x) are equal. Hence the associated
regular distributions are also equal. If a proof by substitution really is required, simply compute,
for any test function
Z
Z +
1 +
H (a x), =
H , (x/a) =
(x/a) dx =
(x) dx = H , .
|a|
|a| 0
0
For a < 0, the reader can check that H (a x) = 1 H (x).
Solution of exercise 8.8. For any test function D we have
x , = , x = [x] (0) = (0) 0 (0) = , .
Similarly, by letting it act on any test function , one shows that x = 2 . Note that this
relation can also be recovered directly by differentiating the relation x = , which yields
+ x = and thus x = 2 .
If f is a locally integrable function of C class, we can define the product f . By the
same method as before, we find that
f (x) (x) = f (0) f (0) .
The assumptions can be weakened by asking only that f be continuous and differentiable
at 0. Without a more complete theory of distributions, the relation above can be taken as a
definition for f (x) (x) in this case.
247

1 sin nx/2 2
, for n = 2, . . . , 7. The
2n sin x/2
limit in the sense of distributions of this sequence is X.
Fig. 8.6 The sequence of functions fn : x 7
Solution of exercise 8.10. We have S = (r R). On the other hand,

(r 2 R 2 ) =
1
1
(r R) =
.
2R
2R S
Solution of exercise 8.11.R By writing the difference of two consecutive terms and integrating by parts, it follows that
1
f (t) dt
1 n
is independent of n. Hence (see the graph in Figure 8.6)
fn (t) dt = 1
n N.
Moreover, fn takes only positive values and converges uniformly to 0 on any interval of the type
]+, 1], with > 0, which shows that
x
f (x) ,
n
2 n

P+
xk
and since fn (x) = k= 2 fn (x) and each term tends to k , it follows that fn X.
n
Of course, it is possible to find simpler sequences of functions that converge weakly to X

in D , for instance, the sequence of functions

x E (x)
fn (x) = n
.
n
Solution of exercise 8.13. Expand the sine as sum of two complex exponentials, rearrange
the terms, and compute the sums of geometric series. The convergence must be justified.
Solution of exercise 8.15. We have E (x, t) = H (x) H (t) and therefore

E
(x, t) = (x) H (t),
x
Distributions II
248
whence
2E
(x, t) = (x) (t).
x t
As for the laplacian, it is equal to E = (x) H (t) + H (x) (t).

i) The force exerted on the truck by the spring is k E (v t) R(t) , so by Newtons law
we can write

d2 R
(t) + k R(t) E (v t) = 0,
m
dt
hence the required formula.
ii) The previous equation can also be written

+ 2 R = 2 E (v t),
and remains valid if E is a distribution. The choice of D+ is justified by the desire to

study a causal system. More precisely, if the excitation (the road) starts at a finite time in
the past, we
to be identically zero before that time. The convolution
want the response

inverse of + 2 in D+ has this property.
iii) In D+ , the solutions are given by R(t) = E (t) H (t) sin t.
In the case of a narrow bump, we have E (t) = (v t) = (t)/v and R(t) =

(/v)H (t) sin t.
For a flat road, we find R(t) = H (t)[1 cos t].
For a wavy road, denoting = v/ we obtain
Zt

sin( t) sin(t)
R(t) = sin(s) sin (t s) ds = H (t)
.
2 2
0
iv) The nitroglycerin will explode if the amplitudes of the vibrations of the truck are too
important. The previous formula indicates a resonance at the speed v = . A very
slow or very fast speed will avoid this resonance.
Solution of problem 3.
Using Theorem 7.51 on page 206, the formula asked in Question 2 follows. To show that ( + k 2 ) G = , use the property
(1/r ) = 4
and the formula for differentiating a product: it follows that

d
(1) e i kr
d
S ( r) +
, S ( r)
,
( r) =
dn
dn
4 r
and then that
( r) =
d
( x) +
dn S
d
dn
S ( x),
(1) e i kk r xk
4 k r xk
Chapter
Hilbert spaces,
Fourier series
9.1
Insufficiency of vector spaces
To indicate how vector spaces and the usual notion of basis (which is called
an algebraic basis) are insufficient, here is a little problem: let E = RN , the
space of sequences of real numbers. Clearly E is an infinite-dimensional vector
space. So, the general theory of linear algebra teaches that1:
THEOREM 9.1 Any vector space has at least one algebraic basis.
The question is then

Can you describe a basis of E ?
(Before turning the page, you should try your skill!)
1
Though the theorem that states this does not give a way of constructing such a basis. Even
worse, the proof uses the axiom of choice or rather an equivalent version, called Zorns lemma
LEMME (Zorn) Let Z be a non-empty ordered set in which any totally ordered subset admits an upper
bound in Z . Then Z contains at least one maximal element.

An elegant statement certainly, but which, for a physicist, looks obviously fishy: not much
that is constructive will come of it!
250
Lets try. One may think that the family

n
(1, 0, 0, 0, . . . ), (0, 1, 0, 0, . . . ), (0, 0, 1, 0, . . . ),
...
(9.1)
of elements of E is a basis. But this is an illusory hope, remembering that a

basis is a free and generating family:
Let I be an arbitrary index set (finite,
countable, or uncountable). The sub-vector space generated
by a family
(x i ) iI of vectors in a vector space E , denoted Vect x i ; i I , is the set of
( finite, by definition) linear combinations of these vectors:
o

def n
P
Vect x i ; i I = x =
a i x i ; I I finite and a i K .
DEFINITION 9.2 (Algebraic basis)
iI
If Vect{x i ; i I } = E , the family (x i ) iI is called a generating family, and

the family (x i ) iI is free if the only (finite) linear combination which is equal
to zero is the one where every coefficient is zero.
The family (x i ) iI is an algebraic basis if it is both free and generating.
This means that any element of E must be expressible as a finite linear

combination of the basis vectors. But it is clear that (1, 1, 1, . . . , 1, . . . ), which
is an element of E , is not a finite linear combination of the family (9.1).
Finite sums, schminite sums! We just need to suppress this condition of
finiteness and permit infinite sums. But what this means needs to be clarified.
More precisely, we must give a meaning to the convergence of a sum
i ei ,
i=0
where the i are scalars and the ei are vectors.

To speak of convergence, the simplest solution is to consider the setting
of normed vector spaces. Indeed, on the space E of real sequences, there is a
natural norm associated to the scalar product
u, v =
X
k=0
uk vk ,
namely,
kuk2 =
X
k=0
|uk |2 .
However, this will not do, since the norm of the vector (1, 1, 1, . . . , 1, . . . )
is infinite. Without additional assumptions, speaking of norms brings other
difficulties, and we are back at the beginning. Well, then, one may wish
to keep the finiteness requirement for the sums, but to restrict the vector
space under consideration by looking only at the space E 0 made of sequences
with bounded support (i.e., sequences where all but finitely many elements
are zero). This time, the family (9.1) is indeed an algebraic basis of E0 , but
Pre-Hilbert spaces
251
another unpleasant phenomenon arises. Consider the sequence
x (1) = (1, 0, 0, 0, 0, . . . ),
x (2) = 1, 12 , 0, 0, 0, . . . ,
x (3) = 1, 12 , 14 , 0, 0, . . . ,
x (4) = 1, 12 , 14 , 18 , 0, . . . ,
1
x (n) = 1, 12 , . . . , 2n1
, 0, . . . ,
...
of elements of E0 .
If we put
x=
1
2k kN
1
= 1, 12 , 14 , 18 , . . . , 2n1
,... ,
it seems obvious that the sequence (x (n) )n converges to x, since
X
(n)

1
1
x x 2 =
=
.
k
4
3 4n1
k=n
Unfortunately, the element x does not belong to E 0 (it does not have bounded
support). One says that E 0 is not complete.
So we can see that the difficulties pile up when we try to extend the notion
of basis to an infinite family of vectors. To stay in a well-behaved setting, we
are led, following Hilbert and Schmidt, to introduce the notion of a Hilbert
space. Before doing this, we will review the general techniques linked to the
scalar product and to projection onto subspaces of finite dimension.
9.2
Pre-Hilbert spaces
DEFINITION 9.3 Let E be a real or complex vector space. A complex scalar
product or a hermitian (scalar) product on E is any application (form)

(x, y) 7 (x| y) such that
(a) (x| y) = ( y|x) for all x, y E ;

(b) x y + y = (x| y) + x y for any C, and x, y, y E ;
(c) (x|x) 0 for all x E ;
(d) (x|x) = 0 if and only if x = 0.

If E is equipped with a hermitian product, it is called a pre-Hilbert space.
252
Property (b) expresses the linearity with respect to the second variable.
The scalar product is thus semilinear (or antilinear) with respect to the first,
which means that

, C, x, x , y, E
x + x y = (x| y) + x y .
The hermitian product is also called definite (property d) and positive (property c).
DEFINITION 9.4 (Norm) Let E be a pre-Hilbert space. The norm (or hermitian norm) of an element x E is defined using the scalar product by
def p
kxk = (x|x).
(Theorem 9.6 on the facing page shows that this is indeed a norm.)
The hermitian product has the following essential property:
2
THEOREM 9.5 (Cauchy-Schwarz inequality) Let E be a pre-Hilbert space. Let
x, y E be two vectors. Then we have

(x| y) kxk k yk ,
with equality if and only if x and y are colinear.
( x |y)
(assuming that (x| y) 6= 0, which is not a restriction since
|( x | y )|
otherwise the inequality is immediate). Then we have = 1 and, for any real number
, the inequalities
def
Proof. Put =
0 (x + y|x + y) = (x|x) + (x| y) + ( y|x) + 2 ( y| y)

= 2 ( y| y) + 2 |(x| y)| + (x|x) .
If k yk2 6= 0 (otherwise the result is also proved), this last quantity is a polynomial
of degree 2 in , positive or zero for any value of R; this is only possible if the
discriminant is negative, namely, if
4 |(x| y)|2 4 ( y| y) (x|x) 0,
p
p
or equivalently if |(x| y)| (x|x) ( y| y), which is the required inequality. Moreover,
there is equality if and only if the polynomial
has

a double root; in that case, taking
to be equal to this root, we obtain x + y = 0, which shows that the vectors x
and y are linearly dependent.
Another importance
inequality, called the Minkowski inequality, shows that
p
the map x 7 (x|x) is indeed a norm (in the sense defined in Appendix A):
2
One should not confuse Hermann Schwarz (18431921) and Laurent Schwartz (1915
2002). For biographical indications, see, respectively, pages 161 and 222.
Pre-Hilbert spaces
253
THEOREM 9.6 (Minkowski inequality) Let E be a pre-Hilbert space. Then we
have
x, y E
kx + yk kxk + k yk ,
with equality if and only if x and y are linearly dependent, and their proportionality
coefficient is positive or zero.
Proof. Use the Cauchy-Schwarz inequality in the expansion of kx + yk2 .
Moreover, we have the following equality:

THEOREM 9.7 (Parallelogram identity) If E is a pre-Hilbert space, then for any
x, y E we have
kx + yk2 + kx yk2 = 2 kxk2 + 2 k yk2 .
This is called the the parallelogram identity (or law) because of the following
analogy with the euclidean norm in the plane: the sum of the squares of the
lengths of the diagonals of a parallelogram is equal to the sum of the squares
of the lengths of the four sides.
x+y
y
x
xy
DEFINITION 9.8 (Orthogonality) Let E be a pre-Hilbert space.
i) Two vectors x and y in E are orthogonal if (x| y) = 0. This is denoted

x y.
ii) Two subsets A and B are orthogonal if, for any a A and any b B,
we have a b. This is denoted A B.
iii) If A is an arbitrary subset of E , the orthogonal of A, denoted A , is
the vector space of elements in E which are orthogonal to all elements
of A:

A = x E ; y A , x y .
DEFINITION 9.9 (Orthonormal system) Let E be a pre-Hilbert space and let

(ei ) iI be an arbitrary family of vectors of E indexed by a set I that may or
not be countable. The family (ei ) iI is

an orthogonal system if (ei e j ) = 0 as soon as i 6= j;

an orthonormal system if (ei e j ) = i j for all i, j.
254
9.2.a
The finite-dimensional case
We assume in this section that E is a pre-Hilbert space of finite dimension

(which is also called a hermitian space in the complex case and a euclidean
space in the real case).
3
Now let (e1 , . . . , en ) be a given orthonormal basis of EP
. If x is any vector
in E , it can be decomposed according to this basis: x = x i ei . Computing
the scalar product of x with one of the basis vectors ek , we obtain, using the
linearity on the right of the scalar product
P
P
n
n
P
n
(ek |x) = ek
xi ei =
x i (ek |ei ) =
x i i,k = xk .
i=1
i=1
i=1
Hence, the coordinates of a vector in an orthonormal basis are simply given

by the scalar products of this vector with each basis vector.
Similarly, if x and y are vectors, their scalar product is
P
P

n
n
n

P
P
n
(x| y) =
xi ei
yj ej =
x i y j (ei e j ) =
xi yi .
| {z } i=1
i=1
j=1
i, j=1
i, j
THEOREM 9.10 (Computations in orthonormal basis) Let (e1 , . . . , en ) be an or-
thonormal basis of E . Then, for any x, y E , we have

x=
(x| y) =
n
P
i=1
n
P
i=1
9.2.b
(ei |x) ei ,
xi yi =
i.e.,
n
P
i=1
(x|ei ) (ei | y) ,
x i = (ei |x) ,
kxk2 =
n
P
i=1
|x i |2 .
Projection on a finite-dimensional subspace
In any pre-Hilbert space, the techniques of the previous section provide a

means to project an arbitrary vector on a subspace V of finite dimension p.
For this, let (e1 , . . . , e p ) be an orthonormal basis of V . Define then
def
xV =
p
X
i=1
(ei |x) ei .
Then x V is an element of V . Moreover, (x x V |ek ) = 0 for k = 1, . . . , p,

hence x x V is orthogonal to V . Finally, for any y V , since y x V is
in V , hence orthogonal to x x V , we can apply Pythagoras theorem and get
kx yk2 = kx x V k2 + k y x V k2 kx x V k2 .
3
To show that such a basis exists is easy: just take an arbitrary basis, then apply the GramSchmidt algorithm to make it orthonormal.
Pre-Hilbert spaces
255
The distance from x to V , which is defined by d (x; V ) = inf yV kx yk, is

therefore attained at a unique point, which is none other than x V .
THEOREM 9.11 (Orthogonal projection theorem) Let E be a pre-Hilbert space
of arbitrary dimension. Let V be a subspace of E of finite dimension. For x E ,

there exists a unique element x V in V such that x x V V , which is given by the
orthogonal projection of x on V , defined by
def
xV =
p
X
i=1
(ei |x) ei .
for any orthonormal basis (e1 , . . . , e p ) of V .

Moreover, x V is the unique vector in V such that d (x; V ) = kx x V k.
x
e2
e1

(e2 x) e2
V
xv

(e1 x) e1
THEOREM 9.12 (Orthogonal complement) Let V be a finite-dimensional sub-
space of E .
i) V is a complementary subspace of V , called the orthogonal complement

of V :
E = V V .
ii) The orthogonal of V is V : V = (V ) .

Proof
i) It is immediate that V V = {0}. Moreover, any vector x is decomposed in

the form x = x V + (x x V ), where the first vector is an element of V and the
second an element of V , which shows that E = V V .
ii) The inclusion V (V ) is clear.

Let now x V . Using the known decomposition E = V V , we can
write x = y + z with y V and z V .Then z is orthogonal to any vector
of V ; in particular, ( y|z) = 0, that is (x z|z) = 0. But x is orthogonal to V
so (x|z). Thus (z|z) = 0: it follows that z = 0, so x = y and hence x V as
desired.
256
Remark 9.13 If V is a subspace of infinite dimension, the previous result may be false. For
instance, consider the pre-Hilbert space E of continuous functions on [0, 1], with the natural
R1
scalar product ( f |g) = 0 f g, and the subspace V of functions f such that f (0) = 0. If
g is a function in V , then it is orthogonal to x 7 x g(x) (which belongs to V ); therefore
2
R1
x g(x) dx = 0 and g = 0. Therefore V = {0}, although {0} is not a complementary
0
subspace of V .
9.2.c Bessel inequality

Let E be a pre-Hilbert space. Consider a countable orthonormal family of E .
For a vector x E , we denote
def
xp =
p
X
n=0
(en |x) en
for any p N,
the orthogonal projection of x on the finite-dimensional subspace V p =

Vect{e0 , . . . , e p }. Since x x p Vect{e0 , . . . , e p }, Pythagoras theorem implies
that

2 2
kxk = x x p + x p ,
2
p
2 X

(x|en ) 2 kxk2 .
and hence x p =
n=0
Since this inequality holds for all p N, we deduce:
THEOREM 9.14 (Bessel inequality) Let E be a pre-Hilbert space. Let (en )nN be
a countable orthonormal system and let x E . Then the series

and we have

X

(en |x) 2 kxk2 .
2
(en |x) converges
n=0
9.3
Hilbert spaces
Recall that a normed vector space E is complete if any Cauchy sequence
in E is convergent in E .
Complete spaces are important in mathematics, because the notion of
convergence is much more natural in such a space. These are therefore the
spaces that a physicist requires intuitively.4
4
Since the formalization of quantum mechanics by Jnos (John) von Neumann (19031967)
in 1932, many physics courses start with the following axiom: The space of states of a quantum
system is a Hilbert space.
Hilbert spaces
257
David Hilbert (18621943), born in Knigsberg in Germany, obtained his first academic position in this town at the age of 22. He
left for Gttingen eleven years later and he ended his career there in
1930. A many-faceted genius, he was interested in the foundations of
mathematics (in particular the problem of axiomatization), but also
in mechanics, general relativity, and number theory, and invented
the notion of completeness, opening the way for the study of what
would later be called Hilbert spaces. His name remains attached to a
list of 23 problems he proposed in 1900 for the attention of mathematicians present and future. Most have been resolved, and some
have been shown to be indecidable.5
DEFINITION 9.15 Let H be a pre-Hilbert space. If H is complete for its

hermitian norm, then it is called a Hilbert space.
9.3.a
Hilbert basis
The spaces we are going to consider from now on will be of infinite dimension
in general. It is harder to work with spaces of infinite dimension since many
convenient properties of finite-dimensional vector spaces (for instance, that
E E ) are not true in general. The interest of Hilbert spaces is to keep
many of these properties; in particular, the notion of basis can be generalized
to infinite-dimensional Hilbert spaces.
DEFINITION 9.16 (Total system) Let H be a Hilbert space. A family (en )nI
of vectors indexed by a set I , countable or not, is a total system if

{ei ; i I } = {0},
or, in other words, if
(x|ei ) = 0, i I
(x = 0).
We now wish to know if it is possible to extend completely the notion

of orthonormal basis of a finite-dimensional Hilbert space, and especially to
recover the formulas of Theorem 9.10 that is, is it possible to write
X
X

(x|en ) 2 ?
x=
(x|en ) en
and
kxk2 =
nN
nN
In what follows, we consider a countable orthonormal system (en )nN , and

define an increasing sequence of subspaces (Vn )nN given by
def
V p = Vect{e0 , e1 , . . . , e p }
5
for any p N.
Such is the case of the continuum hypothesis, a result due to P. Cohen in 1963.
258
Erhard Schmidt (18761959) studied in his native city of Derpt,

in Estonia, then in Berlin (with Schwarz) and in Gttingen.
There, he defended his thesis (supervised by Hilbert) in 1905.
After positions in Zurich, Erlangen, and Breslau, he was hired as
professor at the University of Berlin in 1917. He continued the
ideas of Hilbert concerning integral equations and defined the
notion of a Hilbert space (1905). He also introduced self-adjoint
operators in Hilbert space, was interested in functional analysis
and integral equations, and, in his studies of systems of infinitely
many equations with infinitely many unknowns, introduced geometric notation, which led to the geometric interpretation of
Hilbert spaces.
Also we introduce the sequence of projections of a vector x on the subspaces V p .

DEFINITION 9.17 Let x H . The partial Fourier series of x with respect to
the orthonormal system (en )nN is the sequence of sums S p (x) given by
def
S p (x) =
p
X
n=0
(en |x) en ,
for any p N.
In other words, S p (x) is the orthogonal projection of x on V p . The coefficient

def
cn (x) = (en |x) is called the n-th Fourier coefficient of x.

Does the partial Fourier series converge to x, which is what we would like?
The answer is no in general. However, we have the following result:

PROPOSITION 9.18 For any x H , the sequence Sn (x) nN is a Cauchy sequence for the hermitian norm, and therefore it converges in H .
P
P 2
Proof. From Bessels inequality it follows that the series
|(en |x)|2 =
|cn | has
bounded partial sums; since it is a series with positive terms, it is therefore convergent,
and consequently it is a Cauchy sequence. Hence the quantity
X
2
q
X

q

S p (x) Sq (x) 2 =

c
e
|cn |2
n n =

n= p+1
n= p+1
is the Cauchy remainder of a convergent series of real numbers, and it tends to 0 when
[p, q ]; this shows that the sequence Sn (x) nN is a Cauchy sequence in H which
is, by assumption, complete. So it converges in H .

This is not sufficient to show that the sequence Sn (x) nN converges to x;
as far as we know, it may very well converge to an entirely different vector.
This is the reason for the introduction of the Hilbert basis.
Hilbert spaces
259
DEFINITION 9.19 (Hilbert basis) Let H be a Hilbert space. A

family (ei ) iI
is a Hilbert basis of H if it is a free family such that Vect ei ; i I is

dense6 in H .
An orthonormal Hilbert basis is a Hilbert basis which is also an orthonormal system.
DEFINITION 9.20 (Separable space) A Hilbert space H is separable if there
exists a countable (or finite) Hilbert basis.7
In what follows, all Hilbert spaces considered will be separable.

Take an orthonormal system (en )nN . What does it mean for this system to
be a Hilbert basis? As the next theorem shows, it is a very strong property,
since it provides a unique decomposition of arbitrary vectors as (infinite)
linear combination of the basis vectors (the system is therefore total).
THEOREM 9.21 Let H be a Hilbert space and let (en )nN be an orthonormal system.
The following statements are equivalent:
a) (en )nN is a Hilbert basis, that is, Vect{en ; n N} is dense in H ;

def P
b) for any x H , if we put xn = ni=0 (ei |x) ei , then we have
lim kx xn k = 0;
c) for any x, y H , we have (x| y) =
i=0 (x|e i ) (e i | y);
d) for any x H , we have the Parseval identity kxk2 =
i=0
e) the system (en )nN is a total system: for x H , we have

(ei |x) = 0 for any i N x = 0 .
2
(ei |x) ;
S
Proof. Denote Tn = Vect{e0 , . . . , en } and T = nN Tn .
a) b): by assumption, T is dense in H , hence T = H . Let > 0. There exists
y T such that kx yk . Moreover, since y T , there exists N N such
that y T N , so that, since the sequence (Tn )nN is increasing for inclusion, y Tn
for any n N . In addition the projection of x on Tn achieves the shortest distance
(Pythagoras theorem), hence kx xn k kx yk for
P any n N .
b) c): an easy computation shows that (xn | yn ) = ni=0 (x|e i ) (e i | y) and letting n
go to infinity gives the result.
c) d): immediate.
6
In other words, for a Hilbert basis (e i )iI , any vector in H can be approached, with
arbitrarily high precision, by a finite linear combination of the vectors e i .
7
In certain books, countability is part of the definition of a Hilbert basis.
260
d) e): immediate.
e) a): let x H and denote y = lim xn (from Theorem 9.18, we know that the
n
limit exists). Then (x y|e i ) = 0 for any i N, so that x = y. In particular, the

sequence (xn )nN of elements of T converges to x, and T is therefore dense in H .
So, an orthonormal Hilbert basis is a total orthonormal system and is

something very close to an algebraic basis, since it is a free family and it is
almost generating. For an arbitrary vector x H , the sequence Sn (x) nN
of the partial Fourier series of x with respect to the Hilbert basis (en )nN
converges to x, so that it is possible to write
X
X
X
x=
(en |x) en =
cn (x) en
and
(x| y) =
cn (x) cn ( y). (9.2)
nN
nN
nN
This means that the computational formulas seen in the finite-dimensional

case are equally valid here.
THEOREM 9.22 (Computations in a Hilbert basis) Let (e1 , . . . , en ) be a Hilbert
basis of H . For any x, y H , we have

x=
(x| y) =
n=0
P
n=0
(en |x) en ,
cn (x) cn ( y) =
cn (x) = (en |x) ,
n=0
(x|en ) (en | y) ,
kxk2 =

P
cn (x)2 .
n=0
COROLLARY 9.23 Two elements of a Hilbert space which have the same Fourier
coefficients with respect to a fixed orthonormal Hilbert basis are equal.
Proof. Let x, y H have the same Fourier coefficients with respect to an orthonormal Hilbert basis (en )nN . We then have cn (x y) = cn (x) cn ( y) = 0 for any n N;
hence, according to Property e), it follows that x y = 0.
PROPOSITION 9.24 In a Hilbert space, all Hilbert basis have the same cardinality.
In particular, in a separable Hilbert space, any Hilbert basis is countable.8
COROLLARY 9.25 Every separable Hilbert space is isometric to the space 2 , to the
space L2 (R), and to the space L2 [0, a] (see the next pages for definitions).
Remark 9.26 In quantum mechanics, it is customary to choose a hermitian scalar product

which is linear on the right and semilinear on the left, and we have followed this convention.
However, it should be noted that many mathematics references use the opposite convention. It
is then necessary to write, for instance, cn (x) = (x|en ) instead of (en |x).
Following Dirac, the vector e i is sometimes denoted |e i and the linear form x 7 (e i |x) is
denoted e i |. The effect of the application of a bra x| on a ket | y is conveniently denoted
x| y = (x| y). Then, the property
8
L2 ,
This is important in quantum mechanics. Indeed, since in general the space considered is
which is separable, it follows that the family {| x} xR is not a Hilbert basis.
Hilbert spaces
x| y =
nN
cn (x) cn ( y) =
nN
(x|en ) (en | y) =
261
P
nN
x| en en | y,
which is valid for all x, y H , can be written symbolically as

P
Id =
|en en |
and
x| y = x| Id | y .
nN
Similarly, the formula (9.2) can be interpreted as

P
P
P
|x =
cn (x) en =
(en |x) en =
|en en | x.
nN
nN
nN
| {z }
Id
All of this is detailed in Chapter 14 on page 377.
Two Hilbert spaces will be important in the sequel, the spaces denoted
2 and L2 ; both are extensively used in quantum mechanics as well as in
signal theory. The two following sections give the definitions and the essential
properties of those two spaces. It will be seen that it is possible to translate
from one to the other, and this is precisely the objective of the classical Fourier
series expansion.
9.3.b The 2 space

DEFINITION 9.27 The space of complex-valued sequences (an )nN such that
the series
|an |2 converges is called the 2 space.
THEOREM 9.28 The space 2 , with the hermitian product

def
(a|b) =
n=0
is a Hilbert space.
a n bn ,
Proof. Because of its length and difficulty, the proof is given in Appendix D,
page 602.
A Hilbert basis of 2 is given by the sequence (en )nN , where

def
en = ( 0,
. . , 0} , 1, 0, 0, . . . ).
| .{z
n1 times
(In this is recognized the basis which we were looking for at the beginning
of this chapter!) Since this Hilbert basis is countable, the space 2 is separable.
Remark 9.29 Because N and Z have same cardinality, 2 can also be defined as the space of
families (an )nZ which are square integrable.
262
9.3.c The space L2 [0, a]

Consider the vector space of functions f : [0, a] C, which are square
integrable. One wishes to define a scalar product and norm defined by
Za
Za

f (x)2 dx.
( f , g) 7
f (x) g(x) dx
and
k f k2 =
0
The problem is that if f and g are equal almost everywhere, even without
being equal, then k f gk = 0, so this is not really a norm. The way out
of this difficulty is to identify functions which coincide almost everywhere,
namely, to say9 that if f and g are equal almost everywhere, then they represent
the same object.
DEFINITION 9.30 Let a R+ . The L2 [0, a] space is the space of measur-
able functions defined on [0, a] with values in C which are square integrable,
defined up to equality almost everywhere, with the hermitian scalar product
and the norm given by
Za
Z a
1/2

def
def
f (x)2 dx
f (x) g(x) dx,
and
k f k2 =
.
( f |g) L2 =
0
When speaking of a function in L2 [0, a], the actual meaning implied

will be the equivalence class of this function modulo almost everywhere equality.
THEOREM 9.31 (A Hilbert basis) The system (en )nZ defined by
1
en (x) = p e 2inx/a
a
for any n Z and any x [0, a]
is an orthonormal Hilbert basis of L2 [0, a].

Proof. The proof, which is somewhat delicate, can be found in most real analysis
books, such as [76].
This means that any function in L2 [0, a] can be represented as a sum

(infinite) of periodic functions, in the sense of mean-square convergence of the
partial sums toward the original function. This result is the foundation of the
theory of Fourier analysis of periodic functions (see Section 9.4).
In other words, a function f L2 [0, a] may be represented by a sequence
P
(cn )nZ such that nZ |cn |2 < +, i.e. an element of 2 .
9
Mathematicians speak of taking the quotient by the subspace of functions which are zero
almost everywhere.
Hilbert spaces
9.3.d
263
The L2 (R
R) space
R) space is the space of measurable functions deDEFINITION 9.32 The L2 (R

fined on R, with complex values (up to equality almost everywhere) which are
square integrable, with the norm
sZ
+

def
f (x)2 dx.
k f k2 =
A fundamental result due to Riesz and Fischer, which is not obvious at all,
shows that L2 (R) is complete. It is therefore a Hilbert basis, and it can be
shown to admit a countable Hilbert basis. Thus the Hilbert spaces 2 , L2 [0, a],
and L2 (R) are isometric.
Remark 9.33 Although a Hilbert basis has many advantages, it will also be very interesting to
consider on the space L2 (R) a basis which is not a Hilbert basis, but rather has elements
which are distributions, such as the (noncountable) family x xR , which quantum physicists

denote |x xR , or the family e 2i p x pR , denoted (not without some risk of confusion!)

|p pR . Instead of representing a function f L2 by a discrete sum, it is represented by a
continuous sum, that is, an integral. Such a point of view lacks rigor if a basis in the sense
of distributions is not defined (note, for instance, that x 7 e 2i p x is not even in L2 (R)!).
However, the representation of a function by Fourier transformation
is simply (in another
language) the decomposition of a function in the basis x 7 e 2i p x pR .
Exercise 9.1 Define H to be the space of functions on R, up to equality almost everywhere,
such that ( f | f ) < +, where
Z
2
def
( f |g) =
f (t) g(t) e t dt.
R
Then H with this norm is a Hilbert space. Show that by orthogonalizing the family
(X 0 , . . . , X n , . . .) (using the method due to Erhard Schmidt), one obtains a family of polynomials (H0 , . . . , Hn , . . .), called the Hermite polynomials.
Show (see, for instance, [20, complement B.v]) that
Hn (t) = (1)n e t
dn t 2
e .
dt n
2
The family (H0 , . . . , Hn , . . .) defined by Hn (t) = Hn (t) e t /2 is then a Hilbert basis

of L2 (R) (see, for instance, [47] for a proof). It is used in quantum mechanics, in particular as a Hilbert eigenbasis of the hamiltonian of the harmonic oscillator; it has the interesting
property of being a family of eigenvectors for the Fourier transform.
Remark 9.34 Another Hilbert basis of L 2 (R) is given by the wavelets of C class of Yves
Meyer [37, 54, 65], which generalize the Fourier series expansion on L2 [0, a] (see Section 11.5,
page 321).
264
9.4
Fourier series expansion
It is also possible (and sometimes more convenient) to define the space
L2 [0, a] as the space of functions f : R C which are periodic with period
a and square integrable on [0, a] (since any function defined on [0, a] may be
extended to R by periodicity10).
1
Recall that
en (x) = p e 2inx/a
a
9.4.a
Fourier coefficients of a function
DEFINITION 9.35 Let f L2 [0, a]. For n Z, the n-th Fourier coefficient
of f is the complex number

1
def
cn = p
a
a
0
f (t) e 2int/a dt = (en | f ) .
The sequence of Fourier partial sums is the sequence ( fn )nN defined by

n
n
X
1 X
def
fn (x) = p
ck e 2ik x/a =
ck ek (x).
a k=n
k=n
The Fourier coefficients of the conjugate, the transpose, and the translates
of f are given by the following proposition:
PROPOSITION 9.36 Let f L2 [0, a], and denote by f its conjugate, by f its
transpose and by f , for R, its translate f (x) = f (x ). We have

i) n Z,
cn ( f) = cn ( f ) ;
ii) n Z,
cn ( f) = cn ( f) = cn ( f ) ;
iii) n Z,
cn ( f ) = e 2ina/ cn ( f ).
THEOREM 9.37 If f L2 [0, a] is continuous and is piecewise of C 1 class, then
f L2 [0, a] and has Fourier coefficients given by

cn ( f ) =
2in
c (f)
a n
n Z.
Proof. This is a simple integration by parts.

10
The value at points 0 and a may not coincide, but forgetting this value is of no importance for functions f which are defined up to equality almost everywhere.
265
9.4.b Mean-square convergence

Let f L2 [0, a], so that it can be seen as a function f : R C which is
a-periodic and square integrable on [0, a]. Since L2 [0, a] is a Hilbert space
and the family (en )nN is an orthonormal Hilbert basis of this space (by
Theorem 9.31), the partial Fourier sums of f converge to f in mean-square
(Theorem 9.21):
THEOREM 9.38 (Mean-square convergence) Let f L2 [0, a]. Denote by cn
its Fourier coefficients and by ( fn )nN the sequence of Fourier partial sums of f . Then
the sequence ( fn )nN converges in mean-square to f :
Za

f f n 2 =
f (t) fn (t)2 dt 0.
2
n
DEFINITION 9.39 It is customary to write (abusing notation, since the equality
holds only in the sense of the L2 -norm, and not pointwise) simply that
+
+
X
1 X
f (t) = p
cn e 2int/a =
cn en (t)
(9.3)
a n=
n=
and to call this expression the Fourier series expansion of f .

Remark 9.40 It is important to fully understand why the identity
f =
+
P
n=
does not allow us to write
c n en
P
1 +
f (t) = p
c e 2i nt/a .
a n= n
true
false!
The first equality is proved in the general case of Hilbert spaces. But the Hilbert space which
we use is the space L2 [0, a], where functions are only defined up to almost-everywhere equality,
and where convergence means convergence in mean-square. The space of a-periodic functions,
with the norm defined as the supremum of the absolute values of a function, is not complete.
One cannot hope for a pointwise convergence result purely in the setting of Hilbert spaces.
(This type of convergence will be considered with other tools in Section 9.4.d on page 267.)
The Parseval identity is then written as follows:

THEOREM 9.41 (Parseval identity) Let f L2 [0, a] and let cn denote its Fourier
coefficients. Then the family (|cn |2 )nZ is summable and we have

Za
X

2
f (t)2 dt = f 2 .
|cn | =
2
nZ
Moreover, if f, g L2 [0, a], then ( f g) L1 [0, a] and we have

Za
X
cn ( f ) cn (g) =
f (t) g(t) dt = ( f |g) .
nZ
266
In other words, the representation as Fourier series is an isometry between the

space L2 [0, a] and the space 2 . Sometimes the following real notation is
used (with the same abuse of notation previously mentioned):
+
a
1 X
t
t
f (t) = p0 + p
an cos 2n
+ bn sin 2n
,
2 a
a n=1
a
a
(9.4)
Za
2
t
an = p
f (t) cos 2n
dt = cn + cn ,
a 0
a
Za
2
t
def
bn = p
f (t) sin 2n
dt = i(cn cn ),
a 0
a
def
with
+
Z a

|a0 |2 1 X
2
2
f (t)2 dt.
+
|an | + |bn | =
4
2 n=
0
and
n6=0
Remark 9.42 The Fourier series expansion gives a concrete isometry between the space L 2 [0, a]
and the space 2 . To each square-integrable function, we can associate a sequence (cn )nZ 2 ,
namely, the sequence of its Fourier coefficients.
Conversely, given an arbitrary sequence
P
(cn )nZ 2 , the trigonometric series
cn en converges (in the sense of mean-square convergence) to a function which is square integrable (this is the Riesz-Fischer Theorem). Frdric
Riesz said that the Fourier series expansion was a perpetual round-trip ticket between two
spaces of infinite dimension [52].
9.4.c Fourier series of a function f L1 [0, a]
This business of convergence in L2 norm is all well and good, but it does not
even imply pointwise convergence (and, a fortiori, does not imply uniform
convergence11). Isnt it possible to do better? Can we write (9.3) or (9.4)
rigorously?
We start the discussion of this question by extending the notion of Fourier
series to functions which are not necessarily in L2 [0, a], but are in the larger
space L1 [0, a] of Lebesgue-integrable functions on [0, a].12 Indeed, for f
11
To take an example independent of Fourier series, the sequence (x 7 sinn x)nN converges
to 0 in mean-square, but it does not converge pointwise or uniformly.
12
Notice that, on the contrary, if f L2 [0, a], then f L1 [0, a] since, by the CauchySchwarz inequality we have
Z a
2
Za
Za

2

f (t) dt = | f | , 1 2
f (t)2 dt
1 dt = f 2 a.
0
267
L1 [0, a], it is perfectly possible to define the coefficient

Za
1
def
f (t) e 2int/a dt,
cn = p
a 0
since a function is Lebesgue-integrable if and only if its modulus is integrable.
Hence:
PROPOSITION 9.43 For any function f L1 [0, a], the Fourier coefficients of f
are defined.
It is then possible to construct the sequence ( fn )nN of partial Fourier

sums of f . What is it possible to say about this sequence? Since we have left
the setting of Hilbert spaces, the general results such as Theorem 9.21 can not
be applied. We will subdivide this problem into several questions.
9.4.d Pointwise convergence of the Fourier series
First question: is it true that the coefficients cn tend to 0?

This is not a mathematicians gratuitous question, but rather a crucial
point for the practical use of Fourier series: in an approximate computation,
it is important to be able to truncate the series and compute only with finitely
many terms. A necessary (but not sufficient) condition is therefore that cn 0
as [n ].
The (positive) answer is given by the celebrated Riemann-Lebesgue lemma
(a more general version of which will be given in the next chapter, see Theorem 10.12):
THEOREM 9.44 (Riemann-Lebesgue lemma) Let f L1 [0, a] be an integrable
function. Then we have
a
0
f (t) e 2int/a dt 0.
n
Proof. The proof is in two steps.

Consider first the case where f is of C 1 class; a simple integration by parts implies
the result as follows:
Za
Za

i
ia
f (t) e 2i nt/a dt =
f (t) e 2i nt/a dt +
f (a) f (0) ,
2n 0
2n
0
which, since f and f are both bounded by assumptions, tends to 0 as n tends to

infinity (the integral is bounded).
Then this is extended to all integrable functions by showing that the space of
functions of C 1 class is dense in L1 [0, a], for instance, by applying the theorem of
regularization by convolution (take convolutions of f L1 [0, a] by a gaussian
with
variance 1/p and norm 1, where p is sufficiently large so that f f g p 1 ).
This is technically delicate, but without conceptual difficulty.
268

Joseph Fourier (17681830) studied at the Royal Military School
of Auxerre and, from the age of thirteen, showed an evident interest for mathematics, although he was tempted to become a priest
(he entered a seminary at Saint-Benot-sur-Loire, which he left in
1789). In 1793, he joined a revolutionary committee but, under
the Terror, barely escaped the guillotine (the death of Robespierre
spared him). At the cole Normale Suprieure, Lagrange and
Laplace were his teachers. He obtained a position at the cole
Centrale de Travaux Publics, then was a professor at the cole
Polytechnique. He participated in Napolons Egyptian expedition, was prefect of the Isre, and then studied the theory of heat
propagation, which led him to the expansion of periodic functions.
Second question: does the sequence of partial sums ( fn )nN converge

pointwise to f ?
In the general case, the answer is unfortunately negative. To obtain a
convergence result, additional assumptions of regularity of required on f .
In practice, the result which follows will be the most important.
DEFINITION 9.45 (Regularized function) Let f be a function which admits
a limit on the left and on the right at any point. The regularized function
associated to f is the function
i
1h +
f (x ) + f (x ) .
f : x 7
2
THEOREM 9.46 (Dirichlet) Let f L1 [0, a] and let t0 [0, a]. If both the
limits on the left and on the right f (t0 ) and f (t0+ ) and the limits on the left and
on the right f (t0 ) and f (t0+ ) exist, then the sequence of partial sums fn (t0 ) nN
converges to the regularized value of f at t0 , that is, we have
i
1h +
lim fn (t0 ) =
f (t0 ) + f (t0 ) = f (t0 ).
n
2
Notice that not only the function itself, but also its derivative, must have a
certain regularity.
It is possible to weaken the assumptions of this theorem, and to ask only
that f be of bounded variation [37]. In practice, Dirichlets theorem is sufficient. Sidebar 4 (page 276) gives an illustration of this theorem.
9.4.e
269
Uniform convergence of the Fourier series
Third question: is there uniform convergence of the sequence of partial

sums ( fn )nN to f ?
A positive answer holds in the case of continuous functions which are
piecewise of C 1 class
THEOREM 9.47 Let f : R C be a continuous periodic function, such that the
derivative of f exists except at finitely many points, and moreover f is piecewise

continuous. Then we have:
i) the sequence of partial sums ( fn )nN converges uniformly to f on R;
ii) the Fourier series of f is obtained by differentiating termwise the Fourier series
of f , that is, we have cn ( f ) = 2in cn ( f ) for any n Z;

P+

iii)
n= c n ( f ) < +.
If f is not continuous, uniform convergence may not be true (even if the

Fourier series converges pointwise); this is the Gibbs phenomenon, described
in more detail below.
Fourth question: what link is there between the regularity of f and the
speed of decay of the Fourier coefficients to 0?
(a) We know already that for f L1 [0, a], we have cn 0 (RiemannLebesgue lemma).
P
2
(b) Moreover, if f L2 [0, a], then +
n= |c n | < + (by Hilbert space
theory).
P
(c) Theorem 9.47 shows that if f C 1 [0, a], then +
n= |c n | < +, which
is an even stronger result.
(d) It is also possible to show that if f is of C k1 class on R and piecewise

of C k class, then cn ( f (k) ) = (2in/a)k cn ( f ) and therefore cn = o(1/nk ).
Moreover, the following equivalence holds:
f C [0, a] cn = o(1/nk ) for any k N.
Remark 9.48 It is not sufficient for this last result that f be of C k class on [0, a]: it
is indeed required that the graph of f in the neighborhood of 0 glues nicely with
that in the neighborhood of a.
Remark 9.49 It has been proved (the Carleson-Hunt theorem) that the Fourier series of a
function f which is in L p for some p > 1 converges almost everywhere to f . It is also known
(due to Kolmogorov) that the Fourier series of an integrable function may diverge everywhere.
The settting of L2 functions is the only natural one for the study of Fourier series.
To look for simple (or uniform) convergence is not entirely natural, then. This is why other
expansions, such as those provided by wavelets, are sometimes required.
270
Remark 9.50 Although many physics courses introduce Fourier analysis as a wonderful instru-
ment for the study of sound (music particularly), it should be noted that the expansion in
harmonics of a sound with given pitch (from which Ohm postulated that the timber of the
sound could be obtained) is not sufficient, at least if only the amplitudes |cn | of the harmonics
are considered. The phases of the various harmonics play an important part in the timber
of an instrument. A sound recording played in reverse does not give a recognizable sound,
although the power spectrum is identical [25].
9.4.f
The Gibbs phenomenon
On the figures of Sidebar 4 on page 276, we can see that the partial Fourier
sums, although converging pointwise to the square wave function, do not converge uniformly. Indeed, there always exists a crest close to the discontinuity.
It is possible to show that
this crest get always narrower;
it gets closer to the discontinuity;
its height is constant (roughly 8% of the jump of the function).

An explanation of this phenomenon will be given in the chapter concerning
the Fourier transform, page 311.
EXERCISES
Exercise 9.2 Let f , g : R C be two continuous functions which are 2-periodic. Denote
by h the convolution product of f and g, defined by
Z 2
def 1
h(x) = [ f g](x) =
f (x t) g(t) dt.
2 0
i) Show that h is a 2-periodic continuous function. Compute the Fourier coefficients

of h in terms of those of f and g.
ii) Fix the second argument g. Determine the eigenvalues and eigenvectors of the linear
application f 7 f g.
|sin x| =
for all x R.
8X
sin2 nx
n=1 4n2 1
Exercise 9.4 (Fejer sums) The goal of this exercise is to indicate a way to improve the rep-
resentation of a function by its Fourier series when the conditions of Dirichlets theorem are
not valid.
Let f : R C be 2-periodic and integrable on [0, 2]. Define
Z
n
X
1
Sn f (x) =
ck e i kx ,
where ck =
f (x) e i kx dx.
2
k=n
Exercises
271
The Fejr sum of order n of f is the function

1
n f = (S0 f + + Sn1 f ),
n
that is, the Cesro mean of the n first Fourier partial sums.
i) Show that
Sn f (x) =
where
f (x u) Dn (u) du
and
n f (x) =
f (x u) n (u) du,

1

1 sin n + 2 u
1
sin nu/2 2
and
n (u) =
,
Dn (u) =
2
sin u/2
2n sin u/2
which are respectively called the Dirichlet kernel and the Fejr kernel.
ii) Show that the sequence (n )nN , restricted to [, ], is a Dirac sequence.
iii) Deduce from this Fejrs theorem:13 if f is continuous, the sequence (n f )nN converges
uniformly to f ; if f merely admits right and left limits at any point, then the sequence (n f )nN
converges to the regularized function associated to f .
Remark: the positivity of the Fejr kernel explains the weaker assumptions of Fejrs
Theorem, compared with Theorem 9.47 (where assumptions concerning the derivative
are required).
Exercise 9.5 (Poisson summation formula) Let f be a function of C 2 class on R, such that
lim x 2 f (x) = lim x 2 f (x) = 0.

Show that the Poisson summation formula
+
+ Z +
X
X
f (n) =
e 2i nx f (x) dx
n=
n=
holds.
P
Hint: Define F (x) = nZ f (x + n) and show that F is defined, 1-periodic, and of C 1
class. Compute the Fourier coefficients of F and compare them with the values stated.
This identity will be proved in a more general situation, using the Fourier transforms of
distributions (Theorem 11.22 on page 309).
PROBLEM
Problem 4 (Isoperimetric inequality)
Let u : R C be a function of C 1 class which is periodic of period L. Show that

Z
2
ZL
ZL
L

2

2

u(x)2 dx L
u (x) dx + 1
u(x)
dx
.
2

4
L
0
0
0
Determine the cases where equality holds.
13
Lipt Fejr (18801959), Hungarian mathematician, was professor at the University of
Budapest. He worked in particular on Fourier series, a subject with a bad reputation at the
time because it did not conform to the standards of rigor that Cauchy and Weierstrass had
imposed. For an interesting historical survey of Fejrs theorem and other works of this period,
see [54], which also explores the theory of wavelets, one of the modern developments of
harmonic analysis.
272
Let be a simple closed path of C 1 class in the complex plane, L its length, and A
the area of the region it encloses. Show, using the previous question, that
L2 4A ,
with equality if and only if is a circle.
Hint: Use the complex Green formula (which states that
Z
ZZ
F
dS
F (z, z ) dz = 2i
z
S
S
for any function F continuous with respect to z and z and admitting continous
partial derivatives with respect to x and y) to show that if is parameterized by the
map f : [0, L] C, with f = u + iv and | f |2 = u 2 + v 2 = 1, then the area is equal
to
ZL
A =
u(s) v (s) ds.
0
SOLUTIONS
Solution of exercise 9.2. Using Fubinis theorem, we compute
cn (h) =
1
42
1
=
42
Z
Z
f (x t) g(t) dt e i nx dx
0
2
0
f (x t) e i n(xt) dx
g(t) e i nt dt
= cn ( f ) cn (g).
For an eigenvector f , the coefficients cn ( f ) must therefore satisfy

cn ( f ) cn (g) = cn ( f )
for any n N.
The set of the eigenvalues of f 7 f g is exactly the set of the cn (g)s. Let be an
eigenvalue and let I N be the set of P
indices n such that cn (g) = ; then the eigenfunctions
are the nonzero functions of the kind nI n e i nx .
Solution of exercise 9.3. The function x 7 |sin x| is even and 2-priodic; computing its
Fourier coefficients (only those involving cosines being nonzero) yields
2 X
2
1
1
|sin x| = +
cos 2nx
n=1 2n + 1 2n 1
after a few lines of calculations. Expanding cos 2 = 1 2 sin2 and noticing that

X
1
1
= 1,
2n + 1 2n 1
n=1
it follows that
8 X
sin2 nx
2X
sin2 nx
=
.
n=1 (2n + 1)(2n 1)
n=1 4n2 1
Note that Dirichlets theorem ensures the pointwise convergence of the series (to the value
|sin x|) at any point x R. This convergence is moreover uniform.
|sin x| =
273

i) It suffices to write
n
1 X
e i ku
2 k=n
and to compute this geometric sum. Similarly, we have
Dn (u) =

1
n1
n1
1X
1X
1 Im exp i n + 2 u
n (u) =
D (u) =
,
n k=0 k
n k=0 2
sin u/2
and evaluating the geometric sum gives the stated formula.

R
R
ii) A direct calculation shows that Dn = 1 and so n = 1.
Let > 0 be given. Then for any u [, ], we have sin u/2 sin /2, hence
1
1
u [, ] [, ] ,
0 n (u)
2n sin2 /2
and this upper bound tends to 0 as [n ]; this shows that [, ] [, ]
converges uniformly to 0 and proves therefore that (n )nN (restricted to [, ]) is a
Dirac sequence.
iii) The sequence (n )nN converges weakly to in D , and so the sequence (n f )n
converges uniformly to f by Theorem 8.22 whenever f is continuous; if f has limits
on the right and left, it converges pointwise to the associated regularized function.
The use of Fejr sums resolves the Gibbs phenomenon close to the discontinuities. The
convergence is not uniform, yet it is much more regular, as shows by the figure below, which
represents the Fourier partial sums and the Fejr partial sums of order 25 for the square wave:
0,5
0,5
0,5
0,5
0,5
0,5
0,5
0
0,5
274
One can show [24] that for f L1 [0, 2], the sequence (n f )nN converges to f in L1 ,
that is, we have
Z 2

n f (t) f (t) dt = 0.
lim
n
Solution of exercise 9.5. We compute

up =
=
f (x) e 2i p x dx =
XZ
nZ
XZ
nZ
n+1
f (x) e 2i p x dx
f (x + n) e 2i p(x+n) dx =
0
Z 1X
f (x + n) e 2i p x dx =
0 nZ
XZ
nZ
1
0
f (x + n) e 2i p x dx
F (x) e 2i p x dx = c p (F ),
P
where F (x) =
nZ f (x + p), which is well-defined because of the assumption f (x) =
2
o (1/x ), and where the interversion of the sum and integral are justified by the normal
x
convergence of the series defining F on any compact set.

Moreover, one notes that F is of C 1 class (because the series converges normally on any
compact set) and it can be differentiated termwise because of the condition f (x) = o(1/x 2 ),
which implies the normal convergence of the series of derivatives.
Since F is of C 1 class and obviously is 1-periodic, the Fourier series converges normally. In
particular,
X
X
X
X Z +
F (0) =
f (n) =
c p (F ) =
up =
f (x) e 2i p x dx.
nZ
pZ
pZ
pZ
Denote by cn the Fourier coefficients of the function u.

By Parsevals identity, we have
Z
2
Z
L
+
+

X
X
2
1 L
1

2i nx/L
u(x) dx =
|cn |2 =
u(x)e
dx

L 0
L
0
n=
n=
and hence
Z
2
Z
2
ZL
+
L

L

X

2
1
1

2i
nx/L
u(x) dx =
u(x) dx +
u(x) e
dx

0

L 0
0
n= L
n6=0
Z
2

+
1 L
L2 X
1

=
u(x) dx + 2

L 0
4 n= n2 L
n6=0
Z
2
L

2i nx/L
u (x) e
dx

0

by integration by parts. But, in the last sum, we have n2 1, hence the inequality
Z
2
Z
2
+
+

X
X
1 L
1 L

2i nx/L
2i nx/L
u
(x)
e
dx
u
(x)
e
dx

2

n
L
L
0
0
n=
n=
n6=0
|u (x)|2 dx.
This proves the formula stated.

Since all the terms of the last sum are non-negative, equality can hold only if all
the terms other than n = 1 or n = 1 are zero, which means only if u (x) is a linear
combination of e 2i x/L and e 2i x/L .
275
The path being parameterized by the functions u(s) and v(s), the integration element
dz can be written

dz = u (s) + iv (s) ds.
Take the function F (z, z ) = z to apply the Green formula. The area of the surface
S enclosed by the path = S is
Z
ZL
ZZ
F
dS = z dz =
(u iv)(u + iv ) ds
2iA = 2i
z
O
S
ZL
ZL
ZL
=
(uu + vv ) ds i
vu ds + i
uv ds.
0
The integrand in the first term is the derivative of 21 (u 2 + v 2 ); since u(0) = u(L) and
v(0) = v(L), this first integral is therefore zero.
The second integral can be integrated by parts, which yields
ZL
ZL
vu ds =
uv ds
0
(the boundary term vanishes) and proves that

ZL
A =
u(s) v (s) ds.
0
0, hence the inequality a 2 u 2 + v 2 /a 2 2uv

Z
ZL
a2 L 2
1
A
u ds + 2
v 2 ds.
2 0
2a 0
Now use the inequality proved in , making the assumption that
ZL
u ds = 0
For any a 6= 0, we have (au
and therefore
v /a)2
(translating the path on the real axis if need be, which does not change the area); it
follows that
Z
Z
1 a 2 L2 L 2
1 L 2
A
u
(s)
ds
+
v
(s)
ds
.
2 42 0
a2 0
It suffices now to fix a so that the coefficients in front of the integrals are equal, which
means a 2 = 2/L; with this value, we get
Z
L 1 2
L2
A
(u + v 2 ) ds =
.
4 0
4
Equality holds if au = v /a, which means v = 2u/L. Together with the relation
L
L
u 2 + v 2 = 1, this implies u(s) = 2
cos(2s/L + ) and v(s) = 2
sin(2s/L + ),
with R, which is, indeed, the equation of a circle with circumference L.
Note that this result, which is rather intuitive, is by no means easy to prove!
276
Sidebar 4 (Pointwise convergence and Gibbs phenomenon)
0.5
0.5
1
0.5
0.5
0.5
0.5
0.5
1
0.5
The partial Fourier sums for the square

wave, with 1, . . . , 5 terms:
1
Sn (t) =
0.5
1
0.5
n1
4X
1
sin(2k + 1)t.
k=0 2k + 1
The pointwise convergence at any point

of continuity is visible, but, close to the
discontinuities, it is clear that the convergence is numerically very bad: there
always remains a crest, getting closer
to the discontinuity, with constant amplitude. This is the Gibbs Phenomenon,
first explained by Josiah Gibbs [39, 40].
Chapter
10
Fourier transform
of functions
This chapter introduces an integral transform called the Fourier transform, which
generalizes the Fourier analysis of periodic functions to the case of functions defined
on the whole real axis R.
We start with the definition of the Fourier transform for functions which are
Lebesgue-integrable (elements of L1 ). One problem of the Fourier transform thus
defined is that it does not leave the space L1 stable. It will then be extended
to functions which are square integrable (elements of L2 ), which have a physical
interpretation in terms of energy. The space L2 being stable, any square integrable
function has a Fourier transform which is also square integrable.
10.1
Fourier transform of a function in L1
In this section, we start by defining the Fourier transform of an integrable
function and deriving its main properties.
Although the Fourier transform is a generalization of the notion of a
Fourier series, it is not the L2 or Hilbert space setting which is the simplest.
Therefore we start with functions in L1 (R),1 before generalizing to square
integrable functions.
1
The reason is that, in contrast with functions defined on a finite interval, L2 (R) is not
contained in L1 (R).
Fourier transform of functions
278
10.1.a Definition
DEFINITION 10.1 Let f be a function, real- or complex-valued, depending on
one real variable. The Fourier transform (or spectrum) of f , if it exists, is the
complex-valued function defined for the real variable by
Z
def
fe() =
f (x) e 2i x dx
for all R.
(10.1)
This is denoted symbolically
fe = F [ f ]

fe() = F f (x) .
or
Similarly, the conjugate Fourier transform of a function f is given (when it

exists) by
Z
def
F [ f ] () =
f (x) e 2i x dx.
Remark 10.2 The integral that appears does not always exist. For instance, if we consider
f : x 7 x 2 , the integral
x 2 e 2i x dx
does not exist, for any value of , and f does not have a Fourier transform (in the sense of
functions; it will be seen in the next chapter how to define its Fourier transform in the sense
of distributions).
It is also possible to have to deal with non-integrable functions for which the integral
defining fe() converges as an improper integral. As an example, consider g : x 7 (sin x)/x.
Then g is not integrable but, on the other hand, the limit
ZM
sin x 2i x
e
dx
lim
M +
x
M
exists for all values of . More details about this (not a complete description of what can
happen! this is very difficult to obtain and state) will be given in Section 10.3, which discusses
the extension of the Fourier transform to the space of square integrable functions.
At least if f is integrable on R, it is clear that fe is defined on R. Indeed, recall that for a
function f : R R, the following equivalence holds:
f is Lebesgue-integrable | f | is Lebesgue-integrable,
which shows that if f is integrable, then so is x 7 f (x) e 2i x for any real value of .
THEOREM 10.3 (Fourier transform of an integrable function) Let f be a func-
tion which is Lebesgue-integrable on R. Then the Fourier transform 7 fe() of f is

defined for all R.
Remark 10.4 A number of different conventions regarding the Fourier transform are in use in
various fields of physics and mathematics; in the table on page 612, the most common ones
are presented with the corresponding formula for the inverse Fourier transform.
279
10.1.b Examples
Example 10.5 Consider the rectangle function (already considered in Chapter 7):
def
(x) =
0
1
for |x| > 21 ,
for |x| < 21 .
Then its Fourier transform is a sine cardinal (or cardinal sine):
Z 1
sin
2
def
e =
()
e 2i t dt = sinc() =

12
1
if 6= 0,
if = 0.
Similarly, the Fourier transform of the characteristic function [a,b ] of an interval is
sin (b a) e i (a+b) if 6= 0,
e[a,b ] () = F [a,b ] () =
b a
if = 0.
2
Example 10.6 Consider the gaussian x 7 e x . Then one can show (see Exercice 10.1 on
page 295 or, using the method of residues, Exercise 4.13 on page 126) that
2
2
F e x = e .
Example 10.7 The lorentzian function with parameter a > 0 is defined by
f (x) =
2a
.
a 2 + 42 x 2
Then, by the method of residues, one can show that
2a
= e a|| .
F
a 2 + 42 x 2
The table on page 614 lists some of the most useful Fourier transforms.
10.1.c The L1 space

Let L1 (R) denote temporarily the space of functions on R which are Lebesgueintegrable. Can it be normed by putting
Z

def
k f k = f (t) dt ?
The answer is no. The reason is similar to what happened in the definition
of the space L2 in the previous chapter: if f and g are two functions in L1
which are equal almost everywhere,2 then k f gk = 0, although f 6= g;
this means that the map f 7 k f k is not a norm. To resolve this difficulty,
2
For instance, one can take f = Q , the characteristic function of the set Q, also called the
Dirichlet function, and g the zero function.
280
those functions which are equal almost everywhere are identified. The set of
all functions which are equal to f almost everywhere is the equivalence class
of f. The set of equivalence classes of integrable functions is denoted L1 (R)
or more simply L1 . Thus, when speaking of a function f L1 , what is
meant is in fact any function in the same equivalence class. For instance, the
zero function and the Dirichlet function belong to the same equivalence class,
denoted 0. In order to not complicate matters, there are both said to be equal
to the zero function.
R
In the space L1 , the map f 7 | f | is indeed a norm.
DEFINITION 10.8 The space L1 (R

R ) or L1 is the vector space of integrable
functions, defined up to equality almost everywhere, with the norm

Z

def
k f k1 = f (t) dt.
two functions f and g which coincide almost everywhere. Then

R Consider
R
f = g. Moreover, if is any other function, it follows that f is equal to
g almost everywhere. So applying this result to the case of : x 7 e 2i x ,
it follows that two functions which are equal almost everywhere have the same Fourier
transform.
Hence, we can consider that the Fourier transform is defined on the
space L1 (R) of integrable functions defined up to equality almost everywhere.
10.1.d Elementary properties

THEOREM 10.9 (Properties of e
f ) Let f L1 (R) be an integrable function and
denote by fe its Fourier transform. Then fe is a continuous function on R. Moreover,

it is bounded, and in fact k fek k f k1 .
R) is the vector space of bounded measurDEFINITION 10.10 The space L (R

able functions on R, with the sup norm defined by

def
k f k = sup f (x).
xR
(It is possible to define functions of L up to equality almost everywhere. The

definition of the norm must then be modified accordingly. This refinement
is not essential here).
THEOREM 10.11 (Continuity of the Fourier transform in L1 ) The Fourier trans-
form, defined on L1 (R) and taking values in L (R), is a continuous linear operator, that is
i) (linearity) for any f, g L1 (R) and any , C, we have
F [ f + g] = F [ f ] + F [g] ;
281
ii) (continuity) if a sequence ( fn )nN of integrable functions converges to 0 in the

sense of the L1 norm, then the sequence ( fen )nN tends to 0 in the sense of the
sup norm:
Z

fn (x) dx 0 = sup fen () 0 .
n
Proof of the two preceding theorems. To show the continuity of the function fe,
we apply the theorem of continuity of an integral depending on a parameter (Theorem 3.11 on page 77), using the domination relation

f (x) e 2i x f (x).
x R
Moreover, for any R, we have
Z
Z

fe() = f (x) e 2i x dx f (x) dx = k f k ,
1

which shows that

k fek = sup fe() k f k1 .
R
Since the Fourier transform is obviously linear, this inequality is sufficient to establish
its continuity (see Theorem A.51 on page 583 concerning continuity of a linear map).
The continuity of the Fourier transform is an important fact. It means

that if a sequence ( fn )nN of integrable functions converges to f L1 , in the
sense of the L1 norm, namely,
Z

fn (t) f (t) dt 0,
n
then there will be convergence in the sense of the L norm, that is, uniform
convergence of the sequence ( fen )nN :
L1
cv.u.
fn f = fen fe on R .
Example 10.5 on page 279 shows that the Fourier transform of a function
in L1 is not necessarily integrable itself. So the Fourier transform operator
does not preserve the space L1 , which is somewhat inconvenient, as one often
prefers an honest endomorphism of a vector space to an arbitrary linear map.
(It will be seen in Section 10.3 that it is possible to extend the Fourier transform to the space L2 and that L2 is then stable by the Fourier transform thus
defined.)
On the other hand, one can see that each of the Fourier transforms previously computed has the property that it tends to 0 at infinity. This property
is confirmed in general by the following theorem, which we will not prove
(see, e.g., [76]). It is the analogue of the Riemann-Lebesgue Lemma 9.44 for
Fourier series.
282
THEOREM 10.12 (Riemann-Lebesgue lemma) Let f L1 (R) be an integrable
function and denote by fe its Fourier transform. Then the continuous function fe tends
to 0 at infinity:

lim fe() = 0.
10.1.e Inversion
The question we now ask is: can one inverse the Fourier transform? The
answer is partially yes; moreover, the inverse transform is none other than
the conjugate Fourier transform.
We just saw that, even if f L1 (R), its Fourier transform fe is not itself
always integrable, and so the conjugate Fourier transform of fe may not be
defined. However, in the special case where fe L1 , we have the following
fundamental result.
THEOREM 10.13 (Inversion in L1 (R
R)) Let f L1 (R) be an integrable function
such that the Fourier transform fe of f is itself integrable. Then we have

F [ fe] (x) = f (x)
and also
F [ fe] (x) = f (x)
for almost all x;
at any point x where f is continuous.
In particular, if f is also continuous,3 we have F [ fe] = f .
Proof. We prove the inversion formula only at points of continuity. The proof starts
with the following lemma.
LEMMA 10.14 Let f and g be two integrable functions. Then fe g and f e
g are integrable
and we have
f (t) e
g (t) dt =
fe(t) g(t) dt.
Proof of the lemma. Since fe is bounded and g is integrable, the product

fe g is also integrable and similarly for f e
g . Hence by Fubinis theorem, we
have
Z
Z
Z

f (t) eg (t) dt = f (t)
g(s) e 2i t s ds dt
Z
Z

Z
= g(s)
f (t) e 2i st dt ds = g(s) fe(s) ds,
which proves the lemma.

3
Then we are not looking at it simply up to equality almost everywhere.
283
Now introduce the sequence of functions (hn )nN defined by

hn (x) = exp(2 x 2 /n2 ).
It is a sequence of gaussians which are more and more spread out, and converges
pointwise to the constant function 1. Since fe is integrable and

fe(x) hn (x) e 2i x t fe(x)
for any reals x, t R and any n, Lebesgues dominated convergence theorem proves
that
Z
Z

(10.2)
fe(x) hn (x) e 2i x t dx fe(x) e 2i x t dx = F fe (t)
n
for any t R. The Fourier transform of x 7 hn (x) e 2i x t is, as shown by an immediate

calculation, equal to x 7 e
hn (x t), and the previous lemma yields
Z
Z
fe(x) hn (x) e 2i x t dx = f (x) hen (x t) dx.
(10.3)
is
But, as seen in Example 8.17 on page 231, the sequence of Fourier transforms of hn
n
2 2
e
hn : x 7 p e n x ,
which is a Dirac sequence. It follows that, if t is a point of continuity for f , we have

Z
f (x) hen (x t) dx f (t).
(10.4)
n
Putting equations (10.2), (10.3), and (10.4) together, we obtain, for any point where f is
continuous, that
Z

fe(x) e 2i x t dx = F fe (t) = f (t),
which is the second result stated.
COROLLARY 10.15 If f is integrable and is not equal almost everywhere to a continuous function, then its Fourier transform is not integrable.
Proof. If the Fourier transform fe is integrable, then f coincides almost everywhere
with the inverse Fourier transform of fe, which is continuous.
Remark 10.16 With the conventions we have chosen for the Fourier transform, we have, when
this makes sense, F 1 = F . One could use indifferently F 1 or F . The reader should be
aware that, in other fields, the convention used for the definition of the Fourier transform is
different from the one we use (for instance, in quantum mechanics, in optics, or when dealing
with functions of more than one variable), and the notation is not equivalent then. A table
summarizes the main conventions used and the corresponding inversion formulas (page 612).
Remark 10.17 Theorem 10.13 may be proved more simply using Lemma 10.14 and the fact
that the Fourier transform of the constant function 1 is the Dirac . But this property will
only be established in the next chapter. For this reason we use instead a sequence of functions
(hn )nN converging to 1, while the sequence of their Fourier transforms converges to .
To apply the inversion formula of Theorem 10.13, one needs some information concerning not only f , but also fe (which must be integrable); this may
be inconvenient. In the next proposition, information relative to f suffices
to obtain the inversion formula.
284
PROPOSITION 10.18 Let f be a function of C 2 class such that f , f , and f are
all integrable. Then fe is also integrable; the Fourier inversion formula holds, and we
have
f = F [ fe] .
Example 10.19 Consider again the example of the lorentzian function. We have
f (x) =
2
.
1 + 42 x 2
Then
192 4 x 2 16 2 x 2
,
(1 + 42 x 2 )3

which are both integrable. We deduce from this that F e || = 2/(1 + 42 x 2 ). Moreover,
since all the functions considered are real-valued, it follows by taking the complex conjugate
and exchanging the variables x and that also
f (x) =
16 2 x
(1 + 42 x 2 )2
and

F e |x| =
f (x) =
2
.
1 + 42 2
Hence the inversion formula is also a tool that can be used to compute
new Fourier transforms easily!
10.1.f Extension of the inversion formula

We can now discuss the case where f has a Fourier transform fe which is not
integrable, that is, such that
Z
fe() e 2i x dx
is not defined in the sense of Lebesgue. In some situations, the limit
Z +R
fe() e 2i x d
lim
R+ R
may still be defined. (It is not quite an improper integral, because the bounds
of integration are symmetrical.) For example, the Fourier transform of the
rectangle function is a cardinal sine. The latter is not Lebesgue-integrable,
but the integral
ZR
sin 2i x
def
e
I (x, R) =
d
R
does have a limit as [R ]. In such cases, the following theorem may be
used.
Properties of the Fourier transform
285
THEOREM 10.20 Let f (x) be an integrable function of class C 1 except at finitely
many points of discontinuity, such that f L1 . Then for all x we have

Z +R
i
1h +
lim
fe() e 2i x d =
f (x ) + f (x ) .
R+ R
2

By abuse of notation, this limit will still be denoted F fe (x).
In other words, the inversion formula, in the sense of a limit of integrals

(called the principal value), gives the regularized function associated to f .
This should be compared to Dirichlets theorem on page 268.
In particular, if f is continuous and f and f are integrable, the inversion
formula of Theorem 10.20 gives back the original function f .
Example 10.21 We saw that the Fourier transform of the rectangle function is sinc().
Since the rectangle function has only two discontinuities and since its derivative (in the
sense of functions) is always zero and therefore integrable, it follows that
0 if |x| > 21 ,
Z +R

lim
sinc() e 2i x d = F sinc() (x) = 1 if |x| < 21 ,
R+
R
1
if |x| = 21 .
2
10.2
10.2.a
Transpose and translates
THEOREM 10.22 Let f be an integrable function and a R a real number. We
have
F f (x) = fe(),

F f (x a) = e 2ia fe(),

F f (x) = F f (x) = fe(),

F f (x) e 2i0 x = fe( 0 ).
These results should be compared to the formulas for similar operations on the
Fourier coefficients of the periodic function; see Theorem 9.36 on page 264.
COROLLARY 10.23 In the following table, the properties of f indicated in the first
column are transformed into the properties of fe in the second column, and conversely:
286

f (x)
even
odd
real-valued
purely imaginary
fe()
even
odd
hermitian
antihermitian
Here a function g is a hermitian function if g(x) = g(x) and an antihermitian

function if g(x) = g(x).
10.2.b Dilation
If a dilation is performed (also called a change of scale), that is, if we put
g(x) = f (a x)
x R,
for some
a R ,
the Fourier transform of g can be computed easily:

Z
Z
1 e
dy
2i
x
=
f
,
e
g () = f (a x) e
dx = f ( y) e 2i y/a
|a|
|a|
a
which gives the following theorem.
THEOREM 10.24 Let a R and f L1 . Then we have

1 e
f
.
F f (a x) =
|a|
a
In other words, a dilation in the real world corresponds to a compression

in the Fourier world, and conversely.
10.2.c Derivation
There is a very important relation between the Fourier transform and the
operation of differentiation:
THEOREM 10.25 (Fourier transform of the derivative) Let f L1 be a func-
tion which decays sufficiently fast at infinity so that x 7 x k f (x) belongs also to L1
for k = 0, . . . , n. Then fe can be differentiated n times and we have

F (2ix)k f (x) = fe(k) ()
for k = 1, . . . , n.
Conversely, if f L1 and if f is of C n class and, moreover, its successive

derivatives f (k) are integrable for k = 1, . . . , n, then we have
F f (m) (x) = (2i)m fe()

for k = 1, . . . , n.
287
In particular, it will be useful to remember that

F f (x) = 2i fe()

d e
F 2ix f (x) =
f ().
d
and
Proof. For any fixed x R, the function 7 f (x) e 2i x is of C class, and

its k-th derivative is bounded in modulus by |(2x)k f (x)|, which is integrable by
assumption. Applying the theorem of differentiation under the integral sign, we obtain
Z
Z
i
d e
d h
fe () =
f () =
f (x) e 2i x dx = (2i x) f (x) e 2i x dx,
d
d
and then, by an immediate induction, the first formula stated.
Now for the second part, recall that for any integrable function , we have
Z
ZR
(x) dx = lim
(x) dx.
R+
Since f is assumed to be integrable, an integration by parts yields

ZR
F [ f ] () = lim
f (x) e 2i x dx
R+
R | {z } | {z }
= lim
R+
f (x) e
2i x
iR
(2i) f (x) e
2i x
dx .
As f and f are integrable, f tends to zero at infinity (see Exercise 10.7 on page 296).
The previous formula then shows, by letting R tend to infinity, that
Z
Z
f (x) e 2i x dx = (2i) f (x) e 2i x dx,
which establishes the second formula for k = 1. An induction on k concludes the

proof.
Noticing that if f is integrable and has bounded support then x 7 x k f (x)

is also integrable for any k 0, we get the following corollary:
COROLLARY 10.25.1 If f L1 is a function with bounded support, then its Fourier
transform is of C class.
The formula of Theorem 10.25 on the preceding page also implies a very
important result:
COROLLARY 10.25.2 Let f and g be measurable functions, and let p, q N.
Assume that the derivatives of f , of order 1 to p, exist and are integrable, and that
x 7 x k g(x) is also integrable for k = 0, . . . , q. Then, for any R, we have
Z

fe() |2| p f ( p) (x) dx
and
(q)
e
g ()

|2x|q g(x) dx,
288
These inequalities give very strong bounds: if f is five times differentiable, for
instance, and f (5) is integrable, then fe decays at least as fast as 1/ 5 .
So there exists a link between the regularity of f and the rate of decay of
fe at infinity, and similarly between the rate of decay of f and the regularity
of fe.
10.2.d Rapidly decaying functions
DEFINITION 10.26 A rapidly decaying function is any function f such that

lim x k f (x) = 0
for any k N.
In other words, a rapidly decaying function is one that decays at infinity

faster than any power of x.
2
Example 10.27 The functions x 7 e x and x 7 x 5 (log x) e |x| are rapidly decaying.
This is a useful notion, for instance, because if f is rapidly decaying

and is locally integrable, then x k f (x) is integrable for all k N. From
Theorem 10.25, we deduce:
THEOREM 10.28 If f is an integrable rapidly decaying function, then its Fourier
transform is of C class.
Still using Theorem 10.25, one can show (exercise!) that

THEOREM 10.29 Let f be an integrable function of C class. If f (k) is integrable
for all k N, then fe is a rapidly decaying function.
10.3
The are some functions which are not necessarily integrable, but whose
square is. Such is the sine cardinal function:
def
x 7 sinc(x) =
sin x
x
(extended to R by continuity by putting sinc(0) = 1), which belongs to L2

but not to L1 .
It turns out that in physics, square integrable functions are of paramount
importance and occur frequently:
289
in quantum mechanics, the wave

R function of a particle is a square
integrable function such that |(x)|2 dx = 1;
in optics, the square of the modulus
R of the 2light amplitude represents a
density of energy; the total energy |A(x)| dx is finite;
R
similarly in electricity, or more generally in signal theory, | f (t)|2 dt
is the total energy of a temporal signal t 7 f (t).
It is therefore advisable to extend the definition of the Fourier transform

to the class of square integrable functions. (It will be seen in Chapter 13,
Section 13.6, how to extend this further to functions with infinite total energy
but finite power.) For this purpose, we start by introducing the Schwartz space,
a technical device to reach our goal.4
10.3.a
The space S
DEFINITION 10.30 The Schwartz space, denoted S , is the space of functions
of C class which are rapidly decaying along with all their derivatives.
Example 10.31 Let f D. Then fe is defined and fe S .
Indeed, f has bounded support, so, according to Corollary 10.25.1, fe is infinitely differentiable. Moreover, f can be differentiated p times and f ( p) is integrable for any p N, so
that Corollary 10.25.2 implies that fe decays at least as fast as 1/ p at infinity, for any p N.
Moreover (still by Theorem 10.25), the derivatives of fe are also Fourier transforms of functions
in D, so the same reasoning implies that fe S .
The most important result concerning S is the following:
THEOREM 10.32 The Fourier transform defines a continuous linear operator from
S into itself. In other words, for f S , we have fe S , and if a sequence ( fn )nN

tends to 0 in S then ( fen )nN also tends to 0 in S .
Proof that S is stable. We first show that if f S , we have fe S .

The function f is rapidly decaying, so its Fourier transform fe is of C class. For
all k N, the function f (k) is also rapidly decaying (and integrable). By Theorem 10.25
on page 286, we deduce that fe is rapidly decaying.
There only remains to show that the derivatives of fe are also rapidly decaying. But
by Theorem 10.25, the derivatives of fe are Fourier transforms of functions of the type
x 7 (2ix)k f (x), and it is clear that by definition of S , these are also in S . So
applying the previous argument to those functions shows that fe(k) is rapidly decaying
for any k N.
To speak of continuity of the Fourier transform F [] on S , we must

make precise what the notion of convergence is in S .
4
The Schwartz space will also be useful to define the Fourier transform in the sense of
distributions.
290
DEFINITION 10.33 A sequence ( fn )nN of elements of S converges in S

to 0 if, for any p, q N, we have

(q)
lim sup x p fn (x) = 0.
n xR
This notion of convergence in S is very strong; in particular, it implies the

uniform convergence on R of ( fn )nN together with that of all its derivatives
(k)
fn nN .
To end this discussion of the most important properties of the Fourier
transform on the space S , it is natural to ask what happens for the inversion
formula. We know that if f S , its Fourier transform fe is in S and, since
fe is continuous and integrable, we will have (by Theorem 10.13) the relation
F [F [ f ]] = f . In fact even more is true:
THEOREM 10.34 (Inversion in S ) The Fourier transform is a bicontinuous linear

isomorphism (i.e., bijective and continuous along with its inverse) from S into S .
Its inverse is given by F 1 [] = F [].
10.3.b The Fourier transform in L2

To define the Fourier transform on the space L2 of square integrable functions,
we will use the results of the previous section concerning the Fourier transform
on S , together with the following lemma, which is given without proof.
LEMMA 10.35 The space S is a dense subspace of the space L2 (R).
The next step is to show that, for two functions f and g in S , we have
Z
Z
f (x) g(x) dx =
fe() e
g () d.
This is an easy consequence of Fubinis theorem and some manipulations with

complex conjugation.5 We can guess that a Hilbert space structure is in the
offing!
Then we use the properties
S is dense in L2 ,
L2 is complete,
to show that the Fourier transform defined on S can be extended to a continuous linear operator on L2 . This step, essentially technical, is explained in
Sidebar 5 on page 294.
We end up with an operator Fourier transform, defined for any function
f L2 (R) which satisfies the identities of the next theorem:
5
See, for instance, the proof of Lemma 10.14.
291
THEOREM 10.36 (Parseval-Plancherel) The Fourier transform F [] is an isome-
try on L2 : for any two square integrable functions f and g, the Fourier transforms fe
and e
g also in L2 and we have
Z
Z
f (x) g(x) dx =
fe() e
g () d.
In particular, taking f = g, we have

Z
Z

f (x)2 dx =
or, in other words,

fe()2 d
k f k2 = k fek2 .
THEOREM 10.37 (Inversion in L2 ) For any f L2 , we have

F F [ f] =F F [ f] = f
almost everywhere.
The inverse Fourier transform is thus defined by
F 1 [] = F [].
It is natural to wonder if there is a relation between the Fourier transform
thus defined on L2 (R) (after much wrangling and sleight of hand) and the
original Fourier transform on L1 (R).
PROPOSITION 10.38 The Fourier transform on L1 and that on L2 coincide on the
subspace L1 L2 .
In other words, if f is square integrable on the one hand, but is also integrable
itself on the other hand, then its Fourier transform defined in the section is
indeed given by
Z
F [ f ] : 7
f (x) e 2i x dx.
But then, how does one express the Fourier transform of a square integrable
function which is not integrable? One can show that, in practice, its Fourier
transform in L2 can be defined by the following limit:
ZR
e
f () = lim
f (x) e 2i x dx
for almost all R
R+ R
(which is again a special case of improper integral, defined in the sense of a

principal value).
292
10.4
Fourier transform and convolution
10.4.a Convolution formula
Recall first the definition of the convolution of two functions:
DEFINITION 10.39 Let f and g be two locally integrable functions. Their
convolution, when it exists, is the function denoted h = f g such that

Z +
def
h(x) =
f (t) g(x t) dt.
Fubinis theorem (page 79) provides the proof of the following fundamental theorem, also known as the Faltung theorem (Faltung being the German
word for convolution):
THEOREM 10.40 Let f and g be two functions such that their Fourier transforms
exist and the convolution f g is defined and is integrable. Then we have

F [ f g] = F [ f ] F [g] ,
f g() = fe() e
g ().
Similarly, when those expressions are defined, we have

F [ f g] = F [ f ] F [g] ,
f g() = fe e
g ().
Proof. We consider the special case where f and g are integrable. Then the Fourier
transforms of f and g are defined and f g is indeed integrable. Define (x, t) =
f (t) g(x t) e 2i x . Then
Z Z

Z Z

F [ f g] () =
f (t) g(x t) dt e 2i x dx =
(x, t) dt dx.
Since is integrable on R2 by Fubinis theorem (and the integrability of f and g), we

obtain
Z Z

F [ f g] () =
(x, t) dx dt
Z Z

=
f (t) g(x t) e 2i x dx dt
Z
Z

= f (t)
g(x t) e 2i (xt) dx e 2i t dt
Z
= f (t) eg () e 2i t dt = fe() e
g ().
The second formula is in general more delicate than the first; if we assume that fe
and e
g are both integrable, it follows from the first together with the Fourier inversion
formula.
Fourier transform and convolution
293
Example 10.41 We want to compute without a problem the Fourier transform of the function
defined by
1 + x
(x) = 1 x
for 1 x 0,
for 0 x 1,
for |x| 1.
(x)
1
x
1

First we show that (x) = (x) and then we deduce that

sin 2
F (x) =
.
Note that the Parseval-Plancheral formula gives

Z
Z
Z
(x) dx = (x) dx (x) dx = 1,
which is the expected result.
10.4.b Cases of the convolution formula

One can show, in particular, that the following cases of the convolution formula are valid:
i) if f , g L1 , then
f g() = fe e
g () for any R;
f g() =
ii) if f , g L1 and their Fourier transforms are also in L1 , then
e
f e
g () for any R;
iii) if f and g are in L2 , one can take theinverse

Fourier transform of the
first formula and write f g(t) = F fe e
g (t) for all t R, which is
a good way to compute f g (see Exercise 10.2 on page 295); moreover,
the second formula is valid:
f g() = fe e
g () for all R;
iv) if f L1 and g L2 , then f g(t) = F fe e

g (t) for almost all
t R.
294
Sidebar 5 (Extension of a continuous linear operator) Let E be a normed

vector space and V a dense subspace of E . If F is a complete normed vector
space and : V F is a continuous linear map, then
a) there exists a unique continuous linear extension of , say b : E F (i.e.,

b satisfies (x)
b = (x) for any x V );
b) moreover, the norm of
b and the norm of are equal.
Let x E . We want to find the vector which will be called (x).

b
Using the
fact that V is dense in E , there exists a sequence (xn )nN with values in V which
converges to x in E . To show that the sequence ((xn ))nN converges, we start by
noting that (xn )nN is a Cauchy sequence. Since is linear and continuous, we
have, denoting by |||||| the norm of (see Appendix A, page 583), the inequality

(x p ) (xq ) |||||| x p xq

for any p, q N. The sequence (xn ) nN is therefore also a Cauchy sequence,
and since it takes values in the complete space F , it is convergent. There only
remains to define
(x)
b = lim (xn ).
n
It is easy to check that (x) does not depend on the chosen sequence (xn )nN
converging to x: indeed, if ( y n )nN also converges to x, we have
lim (xn ) lim ( yn ) = lim (xn yn ) = 0
by continuity of .
The map b thus constructed is linear and satisfies

(x)
b = lim (xn ) = lim (xn ) |||||| lim kxn k = |||||| kxk ,
n
showing that ||||||

b ||||||.
To prove the converse inequality, note that for y V , we have (
b y) = ( y)
and therefore

(x)
(x)
b
b
||||||
b = sup
sup
= |||||| .
xE \{0} kxk
xV \{0} kxk
Hence the equality |||
b||| = |||||| holds as stated.
Exercises
295
EXERCISES
Exercise 10.1 (Fourier transform of the gaussian) Using the link between the Fourier trans-
form and derivation, find a differential equation satisfied by the Fourier transform of the
function x 7 exp(x 2 ). Solve this differential equation and deduce from this the Fourier
transform of this function. Recall that
Z
2
e x dx = 1.
Deduce from this the Fourier transform of the centered gaussian with standard deviation :

1
x2
f (x) = p
exp 2 .
2
2
Exercise 10.2 Using the convolution theorem for functions f , g L 2 and the preceding
exercise, compute the convolution product of a centered gaussian f with standard deviation
and a centered gaussian f with standard deviation .
Exercise 10.3 Let C such that Re() > 0, and let n N . Compute the Fourier
transform of the function
x 7 f (x) =
In particular, show that fe() is zero for < 0.
1
.
(2i x )n+1
Exercise 10.4 Compute the following integrals, using the known properties of convolution
and the Parseval-Plancheral formula, and by finding a relation with the functions and
= :
Z +

Z +

Z +
sin x 2
sin x 3
sin x
dx,
dx,
dx.
x
x
x
Exercise 10.5 (sine and cosine Fourier transforms) Let f : R R be an odd integrable func-
tion. Define
Show that
def
fes () = 2
f (t) = 2
f (t) sin(2 t) dt.
+
0
fes () sin(2 t) d.
The function fes thus defined is called the sine Fourier transform of f , and f is the
inverse sine Fourier transform of fes . These transforms can be extended to arbitrary (integrable)
functions, the values of which are only of interest for t > 0 or > 0.
Similarly, show that if f is even and if the cosine Fourier transform of f is defined by
Z +
def
e
fc () = 2
f (t) cos(2 t) dt,
0
then we have
f (t) = 2
fec () cos(2 t) d.
Compute the cosine Fourier transform of x 7 e x , where is a strictly positive real

number.
296
Deduce then the identity
t > 0 > 0
2t
cos 2 t
d =
e
.
2 + 2
2
Exercise 10.6 Given functions a : R R and b : R R, and assuming their Fourier

transforms are defined, solve the integral equation
Z +
y(t) =
y(s) a(t s) ds + b (t).
Exercise 10.7 (proof of Theorem 10.25) Let f L 1 (R) be a function which is differentiable
with a derivative which is Lebesgue-integrable: f L1 (R). Show that lim f (x) = 0.

x
SOLUTIONS
Solution of exercise 10.1. Define f (x) = e x 2 , which is an integrable function, and let fe
be its Fourier transform. From
f (x) + 2x f (x) = 0,
x R
we deduce, evaluating the Fourier transform,

R
2 fe() +
d e
f () = 0,
d
2
which has solutions given by fe() = A e for some constant A to be determined.
R +
Since fe(0) = f (x) dx = 1, it follows that A = 1.
By a simple change of variable, it also follows that

x2
1
F.T.
2 2 2
p exp 2 e 2 ,
2
2
a result which will be very useful in probability theory.

Compare this with the method given in Exercise 4.9 on page 125, which uses the Cauchy
formula in the complex plane.
Solution of exercise 10.2. Using the convolution formula, we get F f f = fe
f
f , and since f and f are both square integrable, their Fourier transforms are also square
integrable. Using the Cauchy-Schwarz inequality, it follows that fe f
f is integrable; there its
conjugate Fourier transform is defined and
2 2 2
2 2 2
f f (t) = F fe fe (t) = F e 2 e 2 (t)
2
2 2
= F e 2( + ) (t) = fp 2 + 2 (t).
The squares of the standard deviations (i.e., the variances) add up when taking the convolution
of gaussians.
Solution of exercise 10.3. We first notice that f is integrable: f L1 . So its Fourier

transform is defined by the integral
fe() =
e 2i x
dx.
(2i x )n+1
297
We compute this using the method of residues. The only pole of the integrand in C is at
x = i/2.
If < 0, we close the contour in the lower half-plane to apply Jordans second lemma,
and no residue appears; therefore the function fe vanishes identically for < 0.
For > 0, we close the contour in the upper half-plane and, by applying the residue
formula, we get fe() = (1)n+1 n e /n!.
To conclude, the Fourier transform of f can be written
n
fe() = (1)n+1 e H ().

n!
Solution of exercise 10.5. An immediate calculation shows that if f is an odd function,

we have fes = (i/2) fe. Moreover, fes is odd; hence the Fourier inversion formula yields
Z +
Z +
Z +
f (t) =
fe() e 2i t d =
fes () sin(2 t) d = 2
fes () sin(2 t) d.
By definition, the cosine Fourier transform of x 7

Z +
2
e t cos(2 t) dt =
0
e x
is given by
2
.
2 + 42 2
The inverse cosine Fourier transform leads to x 7 e |x| , which is the original input function,
after it has been made even. The formula follows.
Solution of exercise 10.6. Assume that the Fourier transform of y also exists. Denote by
Y , A, and B the Fourier transforms of y, a, and b , respectively. Then, by the convolution

theorem, we have
Y () = Y () A() + B()
Y () =
and hence
B()
.
1 A()
So, under the assumption that it exists, the solution is given by the inverse Fourier transform
of the function 7 Y () thus defined.
Solution of exercise 10.7. Since f is integrable, we can write

Z
f (t) dt = lim
x+
Moreover, for any x R we have
f (t) dt.
x
0
f (t) dt = f (x) f (0),
which shows that f admits a limit at +, and similarly at . If those limits were nonzero,
f would certainly not be integrable on R.
Chapter
11
Fourier transform
of distributions
11.1
Definition and properties
The next objective is to define the Fourier transform for distributions.
This is interesting for at least two reasons:
1. it makes it possible to define the Fourier transform of distributions
such as or X;
2. possibly, it provides an extension of the Fourier transform to a larger
class of functions; in particular, functions which are neither in L1 (R),
nor in L2 (R), but which are of constant use in physics, such as the
Heaviside function.
In order to define the Fourier transform of a distribution, we begin, as
is our custom, by restricting to the special case of a regular distribution.
Consider therefore a locally integrable function. But then, unfortunately,
we realize that being locally integrable is not, for a function, a sufficient
condition for the Fourier transform to be defined.
Restricting even more, consider a function f L1 . Being integrable,
it is also locally integrable and has an associated regular distribution, also
Fourier transform of distributions
300
denoted f . Its Fourier transform fe is continuous, hence also locally integrable,

and defines therefore a distribution fe, which acts on a test function by the
formula
Z
Z Z

2i
x
t
e
e
f , = f (t) (t) dt =
f (x) e
dx (t) dt,
which is the same, using Fubinis theorem to exchange the order of integration
(this is allowed because f is integrable and is also), as
Z
Z

Z
fe, = f (x)
e 2i x t (t) dt dx = f (x) (x)
e dx = f ,
e .
This computation justifies the following definition:
DEFINITION 11.1 The Fourier transform of a distribution T is defined by
its action on compactly supported infinitely differentiable test functions ,

which is given by
def

F [T ] , = T , F [] .
Sometimes the notation Te is used instead of F [T ]:

T , = T,
e .
Remark 11.2 It is known that if has compact support
and is
nonzero, F [] does
not have compact support; this means that the quantity T ,

e is not always defined.
Hence not all distribution have an associated Fourier transform. For this reason, the
notion of tempered distribution is introduced.
11.1.a
Tempered distributions
that a rapidly decaying function is a function

DEFINITION 11.3 (S
S ) Recall

f such that lim x k f (x) = 0 for all k N, and that the Schwartz space (of
x
functions which are infinitely differentiable and rapidly decaying along with
all their derivatives) is denoted S .
The topological dual of S (i.e., the space of continuous linear forms on
S ) is denoted S . A tempered distribution is an element of S .
Hence, to say that T is a tempered distribution implies that T is continuous, that is, that
if n 0 in S ,
n
then T , n 0 in C.
n
PROPOSITION 11.4 Any tempered distribution is a classical distribution. In other
words, we have S D .
301
Proof. Note first that D S . Let now T S be a tempered distribution; it acts

on any test function S and hence, a fortiori, on any function D, so T can
be restricted to a linear functional on D.
Moreover, if (n )nN converges to 0 in D, it converges also to 0 in S by the
definition of those two convergences, so T , n tends to 0 in C. This shows that T is
continuous, and hence that T D .
Now let us look for examples of tempered distributions.

Let f be a locally integrable function, which increases at most like a power
function |x|k , where k is some fixed integer, i.e., we have f (x) = O (|x|k ) in
the neighborhood
of . Then for any function which decays rapidly, the
R
integral f exists. The map
Z +
f (x) (x) dx
S 7
is also linear, of course, and continuous which the reader is invited to check.
Therefore, it defines a tempered distribution, also denoted f .
DEFINITION 11.5 (Slowly increasing function) A function f : R C is
slowly increasing if it increases at most like a power of |x| at infinity1.
PROPOSITION 11.6 Any locally integrable slowly increasing function defines a
regular tempered distribution.
Exercise 11.1 Let T be a tempered distribution. Show that for any k N, x k T and T (k)
are also tempered distributions.
Exercise 11.2 Show that the distribution exp(x 4 ) X(x) is not tempered.
(Solution page 326)
11.1.b Fourier transform of tempered distributions

Let T be a tempered distribution. For S , that is, for a C function
which is rapidly decaying, and with all its derivatives also rapidly decaying,
the Fourier transform
F [] is also in the space S (Theorem 10.32). Hence
the quantity T , F [] is defined and, consequently, it is possible to define

the Fourier transform of T .
THEOREM 11.7 Any tempered distribution has a Fourier transform in the sense of
distributions, which is also tempered.
Remark 11.8 Does the Fourier transform of tempered distributions defined in this manner
coincide with the Fourier transform of functions, in the case of a regular distribution?
1
A slowly increasing function, multiplied by a rapidly decaying function, is therefore integrable; the terminology is coherent!
302
Yes indeed. Let f be an integrable function, and denote this time by T ( f ) the associated
regular distribution (to make the distinction clearly). The Fourier transform of g (in the sense
of functions) is a bounded function fe, which therefore defines a regular tempered distribution.
The associated distribution, denoted T ( fe), satisfies therefore for any S
Z
Z

e dx = T ( f ),
e = F T (f ) , .
T ( f ), = fe(x) (x) dx = f (x) (x)
This implies that T ( fe) is indeed equal to the Fourier transform, in the sense of distributions,
of T ( f ).
Similar reasoning holds for f L2 (R) (Exercise!) and shows that the Fourier transform
defined on S coincides with that defined on L2 (R).
The properties of the Fourier transform with respect to scaling, translation,

and differentiation remain valid in the setting of tempered distributions:
THEOREM 11.9 Let T be a tempered distribution. Then

F T = (2i) F [T ] ,

F T (x a) = e 2ia F [T ] ,

F T (n) = (2i)n F [T ] ,

F e 2ia x T = Te ( a).
Proof. Take the first property, for instance. Before starting, remember that the
letters x and , usually used to denote the variable for a function and its Fourier
transform,
respectively,
are in
fact interchangeable.
So, it is perfectly legitimate to write

F [T ] (), () as well as F [T ] (x), (x) .
This being said, let T be a tempered distribution and S . Then we have

F [T ] , = T , F

= T , F 2ix (x)

= F [T ] , 2i x (x) = 2i xF [T ] (x), (x) ,
which is the stated formula (with the variable x replacing ).
The other statements are proved in an identical manner and are left as warm-up
exercises.
Finally, here is the theorem linking convolution and Fourier transform.

THEOREM 11.10 Let S be a tempered distribution and T a distribution with
bounded support. Then the convolution product S T makes sense and is a tempered distribution, and therefore it has a Fourier transform which is given by
F [S T ] = Se Te.
If S is a regular distribution associated to a function of C class, and T a

distribution such that the product S T exists and is a tempered distribution, then its
Fourier transform is given by
F [S T ] = Se Te.
303
Remark 11.11 The preceding formulas need some comment.
Indeed, the product of two

distributions is generally not defined (Jean-Franois Colombeau worked on the possibility of
defining such a product, but his results [23] are far from the scope of this book). However,
it is defined when one of the two distributions is regular and associated with a function of
class C .
In the first formula (F [S T ] = Se Te), the distribution Te is really a regular distribution,
associated with the function f , of class C , defined by

f () = T (t), e 2i t
(the right-hand side is well defined for all because T has bounded support).
11.1.c Examples
The first important example of the Fourier transform of a tempered distribution is that of a constant function.
THEOREM 11.12 (F
F [1] = ) The Fourier transform of the constant function 1 is
the Dirac distribution.
F.T.
1 ().
Proof. Let S be a Schwartz function. We have, by definition and the Fourier
inversion formula,
Z
def

e
1, = 1,
e = ()
e d.
Furthermore, as
e is in S , the inversion formula for the Fourier transform yields
Z
()
e d = F []
e (0) = (0) = , .
This theorem indicates that for a constant function, the signal contains
a single frequency, namely the zero frequency = 0. Similarly, we have:
THEOREM 11.13 The Fourier transforms of the trigonometric functions are given by
for all 0 R.

F e 2i0 x = ( 0 ),
i
1h
( 0 ) + ( + 0 ) ,
F [cos(20 x)] =
2
i
1h
F [sin(20 x)] =
( 0 ) ( + 0 )
2i
This result is quite intuitive: for the functions considered, only the frequency
0 occurs, but it may appear with a sign 1 (see Section 13.4, page 365).
THEOREM 11.14 (F
F [] = 1) The Fourier transform of the Dirac distribution is the
constant function 1 : x 7 1.
F.T.
(x) 1 .
304
The Dirac distribution is so spiked that all frequencies are represented.

Proof. Let S be a Schwartz function. We have
e def

, = , e =
e(0).
But, from the definition of

e, it follows that
Z
(0)
e = (x) e 2i 0x dx = 1, .
| {z }
1(x)
Using Theorem 11.9 for differentiation and translation, we deduce:

THEOREM 11.15 Let m N and a R. Then we have

F = 2i,

F (x a) = e 2ia .

F (m) = (2i)m ,
We now compute the Fourier transform of the Heaviside distribution H

(which is tempered). From H = , we deduce, using the differentiation
e () = 1. But, according to
property of the Fourier transform, that 2i H
Theorems 8.6 on page 229 and 7.28 on page 191, the distribution solutions of
the equation T () = 1 are of the form
() +
1
1
pv ,
2i
with C.
To determine the constant , note that H 12 is a real-valued function which

e 1 , must also be odd. Since
is odd, so that its Fourier transform, namely, H
2
the distribution is even, it follows that = 12 . Denoting by sgn the sign

function, which is given by sgn(x) = 2H (x) 1, we further obtain:
The Fourier transform of
the Heaviside distribution is

1
1
1
F H (x) =
pv
+ .
2i
THEOREM 11.16 (Fourier transform of H and pv 1x )
Similarly,

i
1
1
F H (x) =
pv
+ .
2
Moreover, the distribution pv 1x satisfies
1
F pv
= i sgn ,
x
F [sgn x] =
1
1
pv .
i
305
11.1.d Higher-dimensional Fourier transforms

Consider now distributions defined on Rn . For a function f : Rn R (with
suitable conditions), its Fourier transform is defined on Rn by
Z
Z
def
fe( ) = . . .
f ( x) e 2i x dn x
Rn
for n-uples x and

The vector is a frequency vector. It is often advantageous to replace it
by the wave vector k = 2 . Then the conventions defining the Fourier
transform are
Z
Z
Z
Z
1
def
i
k
x
n
f ( x) e
d x, f ( x) =
...
F ( k) e i k x dn k.
F ( k) = . . .
n
(2)
n
n
R
R
Those definitions provide a way of defining the Fourier transform of a
tempered distribution on Rn , imitating the process in Definition 11.1. The differentiation theorem is easily generalized, and has the following consequence,
for instance:
THEOREM 11.17 (Fourier transform of the laplacian) Let T ( x) be a distribu-
tion in R3 , which admits a Fourier transform. Then the Fourier transform of its
laplacian T ( x) is 4 2 Te ( ).
Consider now the following partial differential equation (coming from particle physicis, plasma physics, or ionic solutions, and explained in Problem 5
on page 325):

m2 f ( x) = 4 ( x).
Taking the Fourier transform of this relation, it follows that
(4 2 m2 ) fe( ) = 4.
This is now an algebraic equation, the solution of which is obvious2:
fe( ) =
m2
4
.
+ 4 2
(11.1)
But is that really true? Not entirely! Since we are looking for distribution solutions, we
must remember that there are many distribution solutions to the equation
(m 2 + 4 2 ) T ( ) = 0:
any distribution which is proportional to (m 2 +4 2 ) will do. But, if the function m 2 +4 2

does not vanish on R3 , on the other hand it does vanish for complex values of . The fact
that fe( ) takes nonzero values for complex values of is more easily interpreted in terms of
a (bilateral) Laplace transform than in terms of a Fourier transform. For instance, a solution
is given by f (x, y, z) = e mx , which does indeed satisfy ( m 2 ) f ( x) = 0, but this f has no
Fourier transform. The solution (11.1) is, in fact, the only one which vanishes at infinity. It is
therefore the one we will keep.
306
There only remains to compute the inverse Fourier transform of this function to obtain a solution to the equation we started with. For this, the method
of residues is still the most efficient. This computation is explained in detail
in the physics problem on page 325. In addition, the whole of Chapter 15 is
devoted to the same subject in various physical cases. The result is
f ( r) =
e mk rk
.
k rk
Letting m to 0, we deduce3 the following result, of great use in electromagnetism:

THEOREM 11.18 (Fourier transform of the Coulomb potential)
In R3 , the Fourier transform of the Coulomb potential r 7 1/ k rk is given by

ZZZ
4
1 i k x 3
F ( k) =
d x= 2.
e
k
R3 k xk
Remark 11.19 The function k 7 1/ k kk2 is locally integrable (in R3 ), and the right-hand side
is therefore a regular distribution.
11.1.e Inversion formula

The inversion formula, known for functions that belong to L1 , has a generalization for tempered distributions. First, let us define the conjugate Fourier
transform of a tempered distribution T S by

S
F [T ] , = T , F [] = T (),
e() .
Then F is the inverse of F , that is to say,
THEOREM 11.20 (Inversion formula in S ) Let T S be a tempered distri-
bution. Its Fourier transform U = F [T ] is also a tempered distribution and

T = F [U ].
3
With no rigorous justification. Here is how to proceed to check this. Start by regularizing
the Coulomb potential by putting, for R > 0,
1/r if r R,
fR ( r) =
1/R if r R.
This regularized Coulomb potential converges to f : r 7 1/r in the sens of distributions (the
verification of this is immediate). Since the Fourier transform is a continuous operator, the
Fourier transform of fR tends to fe as R tends to infinity. After a few lines of computation
(isolating the constant part 1/R), the Fourier transform of fR is computed explicitly and gives
(2)3
4
4
( k) + 2
sin(kR).
R
k
R k3
The first and third terms tend to 0 in the sense of distributions as [R +], as the reader
can easily show.
fR ( k) =
The Dirac comb
307
11.2
The Dirac comb
11.2.a
Recall that the Dirac comb is the distribution defined by

X(x) =
+
X
n=
(x n).
We can then show and this is a fundamental result that

THEOREM 11.21 The Dirac comb is a tempered distribution; hence it has a Fourier
transform in the sense of distributions, which is also equal to the Dirac comb:
F.T.
X(x) X().
This formula can be stated differently; since X is defined as a sum of
Dirac distribiutions, we have
+
+
X
X
F [X] () = F
n () =
e 2in .
n=
n=
Renaming x the free variable (which is perfectly permissible since the

names of the variables are only a matter of choice!) and replacing n by n, as
we can, we obtain therefore:
+
X
e 2inx =
n=
+
X
n=
(x n).
(11.2)
This formula can be interpreted also as the Fourier series expansion of X .

Notice that, in contrast with the case of usual periodic functions, the Fourier
coefficients of this periodic distribution do not converge to 0 (they are all
equal to 1).
Proof of Theorem 11.21. The theorem is equivalent to the formula (11.2), and we will
prove the latter.

Let f (t) = t 2 (t) X(t). This is indeed a function, and its graph is
1
4
f (t) =
1 12
1
2
308
As a periodic function, f can be expanded in Fourier series. Since f is even, only the
cosine coefficients an occur in this expansion, and after integrating by parts twice, we
get
Z 1/2
an = 4
x 2 cos(2nx) dx
0
2

x sin(2nx) 1/2
x
sin(2nx) dx +
(the boundary terms vanish)
n
2n
0
0
Z 1/2

1/2
4
cos(2nx)
8x
=
dx +
cos(2nx)
, (and the integral also)
2
n 0
2n
(2n)
0
= 4
1/2
that is,
(1)n
.
2 n 2
Moreover, differentiating twice the function f , in the sense of distributions, we obtain
first
an =
1
f (t) =
1 21
1
2
and then
f (t) = 2 2
+
X
n=
x n 12 = 2 2 X x 21 .
Hence, differentiating twice the series f (t) =
n=0
2 2X(x 21 ) = 4
+
X
X(x) = 1 + 2
+
X
(1)n e 2i nx
n=1
X(x 12 ) = 1 + 2
and therefore
an cos(2nt) yields
+
X
(1)n e 2i nx
n=1
(1)n e 2i nx e i n = 1 + 2
n=1
(since X is even), which was to be proved.
+
X
n=1
e 2i nx =
+
X
e 2i nx
n=
11.2.b Fourier transform of a periodic function

The case of a periodic function seems at first sight to require some care, since
a 2-periodic function f (for instance), only belongs to L1 (R) or L2 (R) if is
almost everywhere zero. However, any function in L2 [0, a] defines a tempered
distribution, and therefore has a Fourier transform.
The Dirac comb
309
So take f L2 [0, 1], which can be expanded in a Fourier series

f (x) =
+
X
cn e 2inx ,
n=
which converges to f in quadratic mean. By linearity, we easily compute the

Fourier transform of such a series,
+
X
cn ( n).
(11.3)
fe() =
n=
Hence the Fourier coefficients of a function in L2 [0, 1] occur as the

weights of the Dirac distributions which appear in its Fourier transform.
This is what is observed when a periodic signal (for instance, a note which is
held on a musical instrument) is observed on a modern oscilloscope which
computes (via a fast Fourier transform) the spectrum of the signal.
11.2.c Poisson summation formula

Let f be a continuous function and T an arbitrary nonzero real number.
Assume that the following convolution exists:
+
+
x
X
X
1
f (x) X
(x nT ) =
f (x nT )
= f (x)
T
T
n=
n=
(it suffices, for instance, that f decays faster than 1/|x| at infinity). The
f ), or equivalently
Fourier transform of f (x) T1 X (x/T ) is fe() X(T
X n
fe
(T n),
T
nZ
the inverse Fourier transform of which is equal to

1 X e n 2inx/T
e
.
f
T nZ
T
THEOREM 11.22 (Poisson summation formula) Let f be a continuous function

which has a Fourier transform. When the series make sense, we have
+
+
n
X
1 X
f (x nT ) =
fe
e 2inx/T .
T
T
n=
n=
In the special case x = 0 and T = 1, the Poisson summation formula is

simply:
+
+
X
X
f (n) =
fe(n).
n=
n=
310
11.2.d Application to the computation of series

Take, for instance, the problem of computing the sum of the series4
X
1
1
def
,
S=
a 2 + 42 n2 b 2 + 42 n2
n=
with a, b R+ and a 6= b . First, note that if we define
X
1
1
def
S () =
e 2i n
a 2 + 42 n2 b 2 + 42 n2
n=
for all R,
then we have S = S (0) and we can write
S () =
n=
with
fe(n) e 2i n ,
1
1
=e
g () e
h().
42 2 + a 2 42 2 + b 2
To exploit the Poisson summation formula, it is necessary to compute the inverse Fourier
transform of f and hence those of g and h. We know that
2a
F.T.
,
e a|x| 2
a + 42 2
and therefore we deduce that
e a|x|
e b|x|
g(x) =
and
h(x) =
.
2a
2b
From this we can compute f : indeed, from fe = e
g he it follows that f = g h; this convolution
product is only an elementary (tedious) computation (see Exercise 7.5 on page 211) and yields
a|x|

1
e
e b|x|
1 a|x| b|x| 1
e
e
=
.
f (x) =
4ab
2 b 2 a2
a
b
fe() =
Hence we have
S = S (0) =
+
X
n=
fe(n) =
f (n)
n=
b2
1
1 X
1
1 X
e a|n| 2
e b|n| .
2
2
a 2a n=
b a 2b n=
Each series is the sum of two geometric series (one for positive n and one for negative n), and
we derive after some more work that
X
X
1 + e a
1 + e b
e a|n| =
and
e b|n| =
.
a
1e
1 e b
n=
n=
Hence, after rearranging the terms somewhat, we get
1
1
1
1
1
1
1
S= 2
b a 2 a 1 e a
b 1 e b
2ab b + a

1
coth a coth b
=
.
2(b 2 a 2 )
a
b
4
I did not make it up for the mere pleasure of computing a series; it appeared during a
calculation in a finite-temperature quantum field theory problem.
311
Josiah Willard Gibbs (18391903), American physicist, was professor of mathematical physics at Yale University. Gibbs revolutionized the study of thermodynamics in 1873 by a geometric
approach and then, in 1876, by an article concerning the equilibrium properties of mixtures. He had the idea of using diagrams
with temperatureentropy coordinates, where the work during a
cyclic transformation is given by the area of the cycle. It took
a long time for chemists to understand the true breadth of this
paper of 1876, which was written in a mathematical spirit. Gibbs
also worked in pure mathematics, in particular in vector analysis. Finally, his works in statistical mechanics helped provide its
mathematical basis.
11.3
Although the Gibbs phenomenon may be explained purely with the tools
of Fourier series,5 it is easier to take advantage of the Fourier transform of
distributions to clarify things.
Let f be a function, with Fourier transform fe() and denote
Z
f (x) =
fe() e 2i x d
for any > 0. To obtain f , the process is therefore:
decompose f in its spectral components fe();

remove the higher frequencies (those with || > );
reconstruct a function by summing the low-frequencies part of the spectrum, forgetting about the higher frequencies (which characterize the
finer details of the function).
So we see that the function f is a fuzzy version of the original function f . It can be expressed as a convolution, since
sin 2 x
f = F 1 (/2 ) fe() = f (x)
.
x
Now recall that, in the sense of distributions, we have
sin 2 x
lim
= (x),
+
x
5
Using, for instance, the Dirichlet kernel, as in [8].
312
in accordance with the result stated on page 233, which proves that, by forgetting frequencies at a higher and higher level, we will finally recover the
original signal:
lim f = f
(in the sense of distributions).
In the case where f has a discontinuity at 0, we can simply write f =

g +H , where is the jump at the discontinuity and g is now a continuous
functions. Now let us concentrate on what happens to H when put through
the filtering process described above. We compute
Zx
sin 2 y
1
1
sin 2 x
H (x) = H (x)
=
dy = + Si(2 x), (11.4)
x
y
2
where Si is the sine integral, defined (see [2, 42]) by the formula
Zx
sin t
def
Si(x) =
dt.
t
0
If we plot the graph of H1 (x) = 12 + 1 Si(2x), we can see oscillations, with
amplitude around 8% of the unit height, on each side of the discontinuity.6
1
0,8
0,6
0,4
0,2
When the cutoff frequency is increased, formula (11.4) shows that the size
of the oscillations remains constant, but that they are concentrated in the
neighborhood of the point of discontinuity (here, we have = 2, 3, 4, 5 and
the graph has been enlarged in the horizontal direction to make it more
readable):
6
Notice in passing that a low-pass filter in frequencies is not physically possible, since it is
not causal.
313
0,5
What happens now if we consider a function f which is 1-periodic, piecewise continuous, and piecewise of C 1 class? Then f has a Fourier series
expansion
X
cn e 2inx ,
f (x) =
nZ
which corresponds to a Fourier transform

X
fe() =
cn ( n).
nZ
Then the function f , obtained by removing from the spectrum the frequencies || , is none other than
X
f (x) =
cn e 2inx ,
|n|
that is, the partial sum the Fourier series of f of order E [ ] (integral part
of ).
Hence the sequence of partial sums of the Fourier series will exhibit the
same aspect of oscillations with constant amplitude, concentrated closer and
closer to the discontinuities.7
It should be noted that this phenomenon renders Fourier series rather
unreliable in numerical computations, but that it can be avoided by using,
instead of the partial sums of the Fourier series, the Cesro means of those
series, which are called Fejr sums (see Exercise 9.4 on page 270).
In the preceding example with the Heaviside function, the oscillations get narrower in
a continuous manner. Here, since the function, being 1-periodic, has really infinitely many
discontinuities, there appears a superposition of oscillating figures which has the property,
rather unintuitive, of remaining fixed on the intervals [n, n + 1[ and of changing
rapidly when going through the points = n.
314
11.4
Application to physical optics
11.4.a
Link between diaphragm and diffraction figure
We are going to investigate the phenomenon of light diffraction through a

diaphragm. For this purpose, consider an experiment that gives the possibility
of observing the diffraction at infinity (i.e., the diffration corresponding to
the Fraunhofer approximation; see [14] for a detailed study), with a light
source S considered as being purely monochromatic.8
To simplify matters, we consider a system which is invariant under translation along the O y axis. The diaphragm, placed at the focal point of a lens
(with focal length F on the figure), is characterized by a transmittance function f (x) which expresses the ratio between the light amplitude just in front
and behind the diaphragm. Typically, this transmittance function is a characteristic function. For instance, a slit of length (infinitely long) is modeled
by a transmittance function equal to
x
fslit (x) =
.
However, it may happen that one wishes to model an infinitely narrow slit. In
order to still be able to observe a figure, it is required to increase the intensity of
the light arriving in inverse proportion to the width of the slit. This intensity
coefficient is incorporated to the the transmittance, which gives
1 x
= (x).
fslit = lim
0
With this convention, the transmittance function thus becomes a transmittance distribution.
y
F
8
Nonmonochromatic effects will the studied in Section 13.7 on page 371
315
The amplitude which is observed in the direction is, with the Fraunhofer
approximation and linearizing sin , equal to
Z
nt
C f (x) e 2ix/ dx,
up to a multiplicative constant. On a screen at the focal point of a lens with
focal length R, the amplitude is therefore equal to
Z
nt
A( y) = C f (x) e 2i y x/R dx.
Changing the scale for the screen and putting = y/R, which has the
dimension of the inverse of length, we obtain therefore
Z
A() = Cnt f (x) e 2ix dx = fe().
Hence, up to a multiplicative constant (which we will ignore in our computations from now on), the amplitude of the diffraction figure at infinity of a diaphragm
is given by the Fourier transform of its transmittance distribution.
Remark 11.23 From the mathematical point of view, interference and diffraction phenomena
are strictly equivalent.
11.4.b Diaphragm made of infinitely many

infinitely narrow slits
Consider a diaphragm made of infinitely many slits, placed at a fixed distance a from each other, with coordinates equal to na with n ranging over all
integers in Z. The diaphragm is modeled by a Dirac comb distribution
x

x
X
X
(x na) =
a
n = aX
.
f1 (x) =
a
a
nZ
nZ
The diffraction figure produced by the diffraction apparatus thus described is
then (in terms of amplitude) given by the Fourier transform of
x F.T.
f1 (x) = a X
A1 () = a 2 X(a).
a
In other words, infinitely many infinitely narrow slits produce a diffraction

figure consisting of infinitely many infinitely narrow bands. The step between
bands is 1/a.
Thus, using the properties of the Fourier transform, we notice that the larger
the spacing between the slits, the smaller the spacing between the steps of the diffraction
figure.
Of course, the situation described here is never seen in practice, as there
can only be finitely many slits, and they are never infinitely narrow.
316
20
10
-1
-0,5
0,5
1,5
Fig. 11.1 Amplitude of the diffraction figure given by 2N + 1 = 21 infinitely narrow slits.
11.4.c Finite number of infinitely narrow slits

Consider now a diaphragm made of 2N + 1 slits, separated by a fixed distance a. The coordinates of the slits are then N a, . . . , a, 0, a, . . . , N a.
The transmittance of the diaphragm is then modeled by
f2 (x) =
N
X
n=N
(x na).
The diffraction figure is the Fourier transform of f2 , which may be computed by summing directly the Fourier transforms of each Dirac component:
A2 () =
N
X
e 2ina,
n=N
or, summing the partial geometric series,

1 e 2ia(N +1)
2iaN
A2 () = e
1 e 2ia

1
1
e 2ia N + 2 e 2ia N + 2
=
e ia e ia

sin (2N + 1)a
=
.
(11.5)
sin a
This function can be checked to be continuous (by extending by continuity).
It is represented in Figure 11.1; the corresponding intensity is represented in
Figure 11.2. The shape of this curve is, however, not easy to see obviously
by just looking at formula (11.5).
The amplitude A2 may also be obtained using the convolution formulas.
Indeed, the transmittance and the diffraction figure for infinitely many slits
-1,5
-1
-0,5
0,5
317
1,5
Fig. 11.2 Intensity of the diffraction figure produced by 2N + 1 = 21 infinitely narrow

slits. This is the square of the function represented on Figure 11.1.
are linked by
aX
x
F.T.
a 2 X(a),
a
and hence the (2N + 1) slits are modelized, for instance, with the distribution

x
x
f1 (x) = a X
.
a
(2N + 1)a
(We took care that the boundaries of the rectangle function do not coincide with any of the Dirac spikes of the comb, which would have led to an
ambiguity.)
The Fourier transform of a rectangle function with width (2N + 1)a is

sin (2N + 1)a
x
F.T.
(2N + 1) a
,
(2N + 1)a
(2N + 1)a
and therefore
with
F.T.
f2 (x) A2 ()

i
sin
(2N
+
1)a
A2 () = a 2 X(a) (2N + 1) a
(2N + 1)a
+

X
sin (2N + 1)a
n
2
= (2N + 1)a
,

a
(2N + 1)a
n=
h
that is,
F.T.
f2 (x) A2 () = (2N
+ 1)a 2

+
X
sin (2N + 1)(a n)
.
(2N + 1)(a n)
n=
318
So the diffraction figure is given by the superposition of infinitely many

sine cardinal functions centered at k/a with k ranging over Z. But the graph
of the sine cardinal function has the following shape:
namely, a fairly narrow spike surrounded by small oscillations. One can then
easily imagine that the whole diffraction figure is made of a series of spikes,
spaced at a distance 1/a from each other; this is what has been observed in
Figure 11.1.
Notice that the width of the spikes is proportional to 1/(2N + 1), so the
more slits there are in the diaphragm, the narrower the diffraction bands are. The
number of slits may be read off directly from the figure, since it corresponds
to the number of local extrema between two main spikes (the latter being
included in the counting). When, as in the example given here, there is an
odd number of slits, the oscillations from neighboring sine cardinals add up;
this implies in particular that the small diffraction bands between the main
spikes are fairly visible. On the other hand, for an even number of slits, the
oscillations compensate for each other; the secondary bands are then much
less obvious, as can be seen in Figure 11.3, to be compared with Figure 11.1.
Note in passing that we obtained two different formulas for A2 , which
are of course equivalent; this equivalence is, however, difficult to prove using
classical analytic means.
11.4.d Finitely many slits with finite width

If we consider now a diaphragm made up of (2N + 1) slits of width (with
< a so that the slits are physically separated), its transmittance distribution
can be expressed in the form
N
X
x na
f3 (x) =
n=N
N
x
X
=
(x na)
n=N
that is,

x
x
x
f3 (x) = aX
.
a
(2N + 1)a
319
20
10
-0.5
0.5
Fig. 11.3 Diffraction figure of a system with 2N infinitely narrow slits (here 2N =
20). Note that the secondary extremums are much more attenuated than on
Figure 11.1, where 2N + 1 slits occurred.
(Despite the appearance of the last line, this is a regular distribution.) The
transmittance f3 is therefore simply the convolution of the transmittance distribution f2 with a rectangle function with width . The diffraction figure
is therefore given by the product of the Fourier transforms of these two distributions. The Fourier transform of (x/) is a sine cardinal function with
width 1/. In general, in such systems, the width of the slits is very small
compared to the spacing, that is, we have a. This implies that the width
of this sine cardinal function (1/) is much larger than the spacing between
the bands. The amplitude of the diffraction figure
A3 () = A2 ()
sin
is shown in Figure 11.4. One sees clearly there how the intensity of the spikes
decreases. Since the ratio a/ is equal to 4 in this example, the fourth band
from the center is located exactly at the first minimum of the sine cardinal,
and hence vanishes completely.
Remark 11.24 In this last example, three different lengths appear in the problem: the width
of the slits, the spacing a between slits, and the total length (2N + 1)a of the system. These
three lengths are related by
a (2N + 1)a,
and reappear, in the Fourier world, in the form of the characteristic width 1/ of the disappearance of the bands (the largest characteristic length in the diffraction figure), the distance
1/a between the principal bands of diffraction, and the typical width 1/(2N + 1)a of the fine
bands (the shortest), with
1
1
1

.
a
(2N + 1)a
Simply looking at the diffraction figure is therefore sufficient in principle to determine
the relative scales of the system of slits in use.
320
Fig. 11.4 Diffraction figure with 21 slits, with width 4 times smaller than their spacing:
= a/4.
11.4.e Circular lens

Consider now the diffraction by a circular lens with diameter D . It is necessary to use here a two-dimensional Fourier transform. The transmittance
function of the lens is
(
p
0 if x 2 + y 2 > D2 ,
p
f (x, y) =
1 if x 2 + y 2 < D2 .
Since this function

is radially symmetric, it is easier to change variables,
p
def
2
putting r = x + y 2 and writing f (x, y) = F (r) = H (D /2 r).
The amplitude A(, ) of the diffraction figure is then given by
ZZ
A(, ) =
f (x, y) e 2i(x+ y) dx dy,
which is also, substituting (x, y) 7 (r, ) and defining by =

equal to
Z D /2
Z
A() = 2
r F (r)
e 2ir cos d dr.
DEFINITION 11.25 The Bessel function of order 0 is defined by
z C
1
J0 (z) =
2
def
e i z cos d
The Bessel function of order 1 is defined by

Z
def 1
z C
J1 (z) =
e i z cos cos d
i
Rx
These functions satisfy 0 y J0 ( y) dy = x J1 (x).
p
2 + 2 ,
Limitations of Fourier analysis and wavelets
321
2
0
Fig. 11.5 The function x 7 2 J1 (x)/x. The first zero is located at a point x 1.220,
the second at x 2.233. The distance between successive zeros decreases.
The amplitude A() can therefore be expressed as

Z D /2
Z D
1
D
A() = 2
r J0 (2r) dr =
y J0 ( y) dy =
J (D )
2
2 0
2 1
0

2 J1 (D )
D 2
=S
,
with S =
= surface of the lens.
D
4
The graph of the function x 7 2 J1 (x)/x is represented in Figure 11.5.

This function has a smallest positive zero at x 1.22, which corresponds
to an angle equal to 1.22 /D ; this is the origin of the numerical factor
in the Rayleigh criterion for the resolution of an optical instrument. The
diffraction figure has the shape presented in Figure 11.6.
11.5
Fourier analysis can be used to extract from a signal t 7 f (t) its composing frequencies. However, it is not suitable for the analysis of a musical or
vocal signal, for instance.
The reason is quite simple. Assume that we are trying to determine the
various frequencies that occur in the performance of a Schubert lied. Then
we have to integrate the sound signal from t = to t = +, which
is already somewhat difficult (the recording must have a beginning and an
ending if we dont want to spend a fortune in magnetic tape). If, moreover,
knowing fe we wish to reconstruct f , we must again integrate over ], +[.
Such an integration is usually numerical, and apart from the approximate
knowledge of fe, the integration range must be limited to a finite interval,
say [0 , 0 ]. But, even if fe decays rapidly, forgetting the tail of fe may
spectacularly affect the reconstitution of the original signal at these instants
when the latter changes very suddenly, which is the case during the attack
322
Fig. 11.6 The diffraction figure of a circular lens. Other larger rings appear in practice,
but their intensity is much smaller.
of a note, or when consonants are sung or spoken. Thus, to recreate the attack
of a note (a crucial moment to recognize the timber of an instrument), it is
necessary to compute the function fe() with high precision for large values
of , and hence to know precisely the attack of all the notes in the partition.
When the original signal is recreated, their musical characteristics and their
temporal localization may be altered.
This is why some researchers (for instance, the geophysicist Jean Morlet,
Alex Grossman in Marseilles, Yves Meyer at the cole Polytechnique, Pierre
Gilles Lemari-Rieusset in Orsay and then vry, and many others) have developed a different approach, based on a time-frequency analysis.

We know that the modulus of the Fourier spectrum fe() of a signal
provides excellent information on the frequency aspects, but little information
on the temporal aspect.9 However, to encode a musical signal, the customary
way is to proceed as follows:
x
On this staff, we can read:

vertically, the frequency information (the height of the notes);
horizontally, the timing information (the rhythm, or even indications

concerning attack).
A given signal f (t) and the same translated in time f (t a) have the same power spectrum,
since only the phase changes between the two spectra.
323
The musical signal is therefore encoded in matters of both time and frequency.
Yves Meyer has constructed a wavelet , that is, a function of C class,
all moments of which are zero, that produces an orthonormal basis (n p )n,p
of L2 , where
def
n p (t) = 2n/2 (2n t p)
n, p Z.
The larger n Z is, the narrower the wavelet n p is; when n , on
the other hand, the wavelet is extremely wide. Moreover, the parameter p
characterizes the average position of the wavelet. The decomposition in the
wavelet basis is given by
f =
+
X
+
X
f n p n p
n= p=
f L2 (R),
in the sense of convergence in L2 . Similarly, a formula analogous to the

Parseval identity holds:
X 2 Z +

f n p =
f (t)2 dt.
(n,p)Z2
Whereas the basis for Fourier expansion was a complex exponential,

which is extremely well localized in terms of frequency10 but not localized at
all in time, for wavelets, the time localization is more precise when n is large,
and the frequency localization is wider when n is large. Thus, both time and
frequency information are available. This makes it possible to describe the
attack of a note at time t = t0 by concentrating on the wavelet coefficients
localized around t0 ; these, even at the highest frequencies, may be computed
without any knowledge of the signal at other instants.
Research around wavelets is still developing constantly. The interested
reader is invited to consult the book by Kahane and Lemari-Rieusset [54]
for a detailed study, or Gasquet [37] for a very clear short introduction, or
Mallat [65] for applications to the analysis of signals. The book by Stein and
Shakarchi [85] should also be of much interest.
10
A single frequency occurs.
324
EXERCISES
Exercise 11.3 Compute (x/a) sin x using the Fourier transform.
Exercise 11.4 Compute the Fourier transform, in the sense of distributions, of x 7 |x|.
Exercise 11.5 Compute the Fourier transform of H (x) x n for any integer n. (One may use
the distributions fp(1/x k ), introduced in Exercise 8.6 on page 241.)

Exercise 11.6 Compute the sum
def
S=
X
nN
a2
1
1
,
2
2
2
+ 4 n b 42 n2
with a, b
The result of Exercise 8.13 on page 243 may be freely used.
R+ .
Exercise 11.7 Recall that the Bessel function of order 0 is defined by

def
J0 (x) =
1
2
e i x cos() d
for any x R.
i) Show that this definition is independent of . Deduce from this that J0 takes real
values.
ii) In the plane R2 , let f (x, y) be a radial function or distribution, that is, such that there
exists : R+ C (or a distribution with support in R+ ) with f (x, y) = (r ), where
p
def
r = x 2 + y 2 . Assume moreover that f has a Fourier transform. Show that its Fourier
def p
transform is also radial: fe(u, v) = () with = u 2 + v 2 .
Prove the formula giving as a function of , using the Bessel function J0 .
H. T.
The function is called the Hankel transform of , denoted (r ) ().

H. T.
H. T.
iii) Show that for continuous functions and , if , then also (i.e.,
the Hankel transform is its own inverse).
iv) Show that
d2 1 d H. T.
+
42 2 ().
dr 2
r dr
H. T.
H. T.
v) Show that if 1 1 and 2 2 , then we have the relation

Z
Z
r 1 (r ) 2 (r ) dr =
1 () 2 () d.
0
Physical optics
Exercise 11.8 Describe the diffraction figure formed by a system made of infinitely many
regularly spaced slits of finite width (the spacing between slits is a).
Exercise 11.9 Compute, in two dimensions, the diffraction figure of a diaphragm made of
two circular holes of diameter , placed at a distance a > 2.
Exercises
325
PROBLEM
Problem 5 (Debye electrostatic screen)
Presentation of the problem

We consider a plasma with multiple components, or in other words, a system of particles of different types I , where I is a finite set, each type being characterized
by a density , a mass m , and a charge e . (Examples of such systems are an ionic
solution, the stellar matter, the solar wind or solar crown.) It is assumed that the system
is initially neutral, that is, that
X
e = 0.
The system is considered in nonrelativistic terms, and it is assumed that the interactions
between particles are restricted to Coulomb interactions. The particles are also assumed
to be classical (not ruled by quantum mechanics) and point-like. The system may be
described by the following hamiltonian:
H =
X p2
1 X ei e j
i
+
2m
2 i, j ri j
i
i
i6= j
with

def
ri j = r j r i .
The indices i and j ranges successively over all particles in the system.
Now an exterior charge q0 , located at r = 0, is added to the system. We wish to know
the repartition of the other charges where thermodynamic equilibrium is reached.
Solution of the problem
i) Let ( r ) denote the distribution of density of particles of type , at equilibrium,
in the presence of q0 . Express the distribution q( r ) in terms of the ( r).
ii) Denoting by ( r) the electrostatic potential, write the two equations relating
( r) and ( r) the Poisson equation for electrostatics and the Boltzmann
equation for thermodynamics.
iii) Linearize the Boltzmann equation and deduce an equation for ( r).
iv) Passing to Fourier transforms, compute the algebraic expression of (
e k), the
Fourier transform of ( r).
v) Using complex integration and the method of residues, find the formula for the
potential ( r), which is called the Debye potential [26].
vi) Show that, at short distances, this potential is indeed of Coulomb type.
vii) Compute the total charge of the system and comment.
326
SOLUTIONS
Solution of exercise 11.2 on page 301. It must be shown that T (x) = exp(x 4 ) X(x) can-
not act on certain functions of the Schwartz space. It is easy to check that the gaussian
g : x 7 exp(x 2 ) decays rapidly as do all its derivatives, and therefore belongs to S , whereas
the evaluation of T , g leads to an infinite result.
Strictly speaking, we must exclude also the possibility that there exists another definition of
T , , for a Schwartz function, which makes it a tempered distribution and coincides with
R
T (x) (x) dx for with bounded support; but that is easy by taking a sequence n of test
functions with compact support that converges in S to g(x). Then on the one hand
T , n T , g
by the assumed continuity of T as tempered distribution, and on the other hand

Z
T , n = T (x) n (x) dx +,
a contradiction.
Solution of exercise 11.3. The Fourier transform of the function given is
1
1
a 1
1
1
sin a 1

+
= 2 sin

+
2i
2
2
2 2i
2
2
(using the formula f (x) (x) = f (0) (x)), which, in inverse Fourier transform, gives
2 sin(a/2) sin x.
Of course a direct computation leads to the same result.
Solution of exercise 11.4. Notice first that f : x 7 |x| is the product of x 7 x and x 7
sgn x. The Fourier transform of x 7 x is, by the differentiation theorem, equal to i /2.
F.T.
Moreover, we know that sgn pv 1/i. Using the convolution theorem yields
F.T.
x sgn x
i 1
1
1 d
1
1
1
pv =
pv = 2 fp 2 .
2 i
22 d
Solution of exercise 11.5. The Fourier transforme of the constant function 1 is . The
formula F [(2ix)n f (x)] = f (n) (x) gives
F.T.
x n
(n)
.
(2i)n
and
The convolution formula then yields

F.T.
H (x) x n
F.T.
H (x)
1
1
+
pv
.
2
2i

(n)
1
dn
1
pv
.
n
n+1
n
2 (2i)
(2i) d
Finally, the distribution fp(1/x n+1 ) is linked to the n-th derivative of pv(1/) by a multiplicative
constant:

dn
1
1
n
pv
=
(1)
n!
fp
,
d n
n+1
from which we derive
(n)
n!
1
F.T.
H (x) x n
+
fp
.
2 (2i)n
(2i)n+1
n+1
327
Solution of exercise 11.6. Using the technique described in the text on page 310, we first
compute that
2a
2b
e a|x| + 2
sin b |x| ,
a2 + b 2
a + b2
and using the formula of exercise 8.13 on page 243 we get
e a|x| sin b |x| =
S=
1
1
sin b
coth a +
.
2
2
2
+b )
2b (a + b ) 1 cos b
2a(a 2

i) A simple change of variable and the periodicity of the cosine function imply that the
integral does not depend on R. Moreover, putting = /2, a sine appears, and as
this is an odd function, it follows that the integral is real.
ii) Passing in polar coordinates (x, y) 7 (r, ), we get
ZZ
fe(u, v) = () =
e 2i (ux+v y) dx dy
=
R2
+
(r )
() = 2
or
Z

e 2i r cos d r dr,
r (r ) J0 (2r ) dr.
0
H. T.
This will be denoted, as indicated, by (r ) ().
iii) This is simply the translation of the Fourier inversion formula in this context.
iv) In the formula
d2 1 d
+
dr 2
r dr
we recognize the expression for the laplacian of f in polar coordinates, since there is
no dependency on . The Fourier transform of f is equal to 42 2 fe(u, v) with
= (u, v) and 2 = 2 , hence the formula follows.
v) The Parseval-Plancherel identity

Z
Z
f1 (x, y) f2 (x, y) dx dy = fe1 (u, v) fe2 (u, v) du dv,
gives, in polar coordinates, the required identity:
Z
Z
2
r 1 (r ) 2 (r ) dr = 2
0
1 () 2 () d.
Solution of exercise 11.9. As seen on page 320, the amplitude of the diffraction figure is
given radially by the Hankel transform of the radial transmittance function, which is given here
by
0 if r > ,
f (x, y) = f (r ) =
1 if r < ,
p
2
2
where r = x + y .
The amplitude is therefore
p
A(, ) = A() =
J1 (2)
with
= 2 + 2 .
328
a
2
a2
Fig. 11.7 The diffraction figure created by two circular holes. The spacing between the
two lenses is here taken to be 4 times their diameter.
The transmittance function of two holes is
h

a i
a
+ x +
.
g(x, y) = f (x, y) x
2
2
The diffraction figure due to the two holes is obtained by multiplying the diffraction amplitude
due to a single hole by the Fourier transform of the distribution (x a/2) + (x + a/2),
which is equal to 2 cos(a).
We obtain then the result sketched in Figure 11.7.
The total density of charge is

X
q( r) =
e ( r) + q0 ( r).
Denote by ( r) the effective electrostatic potential created by the distribution of charge q( r )

(including therefore both the exterior charge and the induced polarization cloud). The quantities ( r) and q( r ) are related by two equations. The first, an electrostatic relation, is the
Poisson equation, which is
( r) = 4 q( r ) ;
()
in Gauss units. The second, a thermodynamic relation, is the Boltzmann equation
( r) = e e ( r ) ,
()
where the constant value of is justifed by the fact that, for r , we have ( r) 0 in
the suitable gauge and ( r) .
In the limit of weak couplings (or low densities), we can linearize equation () and obtain

( r) = 1 e ( r) ,
which, inserted into (), gives
( 2 ) ( r) = 4q0 ( r)
def
with 2 = 4
linearized Poisson-Boltzmann
P
taking into account that the system is neutral ( e = 0).
e2 ,
329
We will now try to solve this equation. First, assuming that ( r) has a Fourier transform,
we obtain the conjugate equation
( k 2 2 ) (
e k) = 4q0
with the conventions

Z
def
(
e k) = e i k r ( r) d3 r
and
( r) =
We deduce that11
1
(2)3
e i k r (
e k) d3 k.
4q0
.
k 2 + 2
There only remains to compute the inverse Fourier transform
Z
4q0 i k r d3 r
.
( r) =
e
(2)3
k 2 + 2
For this purpose, we use spherical coordinates on the space of vectors k, with the polar
axis oriented in the direction of r , and with the notation k r = kr cos , where k > 0 and
r > 0. Then
Z 2
Z Z +
4q0 2
dk
( r) =
d
d
k sin e i kr cos
,
2 + 2
3
k
(2)
0
0
0
or
Z1
Z +
82 q0 2 i kr cos dk
( r) =
d(cos )
k e
k 2 + 2
(2)3
1
0
Z +
Z

q0 k 2
1 i kr
2q0 + k sin kr
i kr
=
e
e
dk
=
dk.
k 2 + 2 ikr
r 0
k 2 + 2
0
(
e k) =
We must still evaluate this integral relative to the variable k. For this, we put
Z +

Z +
1
k e i kr
k sin kr
def
dk
=
Im
dk
,
I ( r) =
2
2
k 2 + 2
2
k +
0
and we use the method of residues by extending the variable k to the complex plane. The poles
def
of the meromorphic function f (z) = z e i zr /(z 2 + 2 ) are located at z1 = i and z2 = i.
Since r > 0, we can integrate on the contour described in the following figure:
z2
z1
and, after applying Jordans second lemma 4.85, as justified by the fact that k/(k 2 + 2 ) tends
to 0 when [k ], and after taking the limit where the radius in the contour goes to infinity,
we derive
Z +
k e i kr
ik e r
dk = 2i Res ( f ; z1 ) = 2i
= i e r ,
2
2
2i
k +
11
Being careful of the fact that, in the general case, the equation x T (x) = a has solutions
given T (x) = a pv(1/x) + b , with b C. Here, (
e k) is defined up to the addition of the
Fourier transform of an harmonic functions.
330
hence
I =
and
r
e
2
( r) =
q0 r
e
r
Debye potential.
Using the Boltzmann equation, each of the functions ( r) can be computed now. Note
that the Debye electrostatic screen occurs at a length scale of 1 , which tends to infinity in the
limit of weak couplings (small densities). For distances r , the potential is ( r) q0 /r ,
and we recover the Coulomb potential.
The total charge present around q0 is
ZZZ X
Q=
e ( r ) d3 r
ZZZ X

e 1 e ( r) d3 r
ZZZ X
= q0
e2 ( r) d3 r
by initial neutrality,
ZZZ X
e2
e r 3
d r = q0 .
r
Thus is exactly compensates the additional charge q0 . In other words, the system has managed
to become neutral again.12 The Debye screen is said to be a total screen.
12
This may seem surprising; whichever charge is added to the system, it remains neutral.
This is a strange property of an infinite system, which may bring in charges from infinity, and
make them disappear there. This type of paradox also happens in the famous infinite hotels,
all rooms of which are occupied, but in which, nevertheless, a new customer can always be
accommodated: it suffices that the person in Room 1 move to Room 2, the person in Room 2
move to Room 3, and more generally, the person in Room n move to Room n +1. All previous
occupants have a new room, and Room 1 becomes available for the new customer.
Chapter
12
The Laplace transform

The Laplace transform is an integral transformation, which is a kind of generalization of the Fourier transform. It has a double interest:
1. it avoids using distributions in some cases where a function does not have a
Fourier transform in the sense of functions;
2. it can be used to solve problems described by differential equations taking
initial conditions into account, that is, it can attack the Cauchy problem
associated to the system. An example will be given to the dynamics of the
electromagnetic field without sources.
12.1
Definition and integrability
In this chapter, we are interested in an integral transformation operating
on functions f which vanish for negative values of the variable: f (t) = 0
for all t < 0 (no continuity at 0 is imposed). An example is the function
t 7 H (t) cos t.
In the literature about the Laplace transform, the factor H (t) is frequently
omitted; we will follow this custom, except where some ambiguity may result,
and we will therefore speak of the function t 7 cos t, tacitly assuming that
this definition is restricted to positive values of t.
DEFINITION 12.1 A causal function is a function t 7 f (t) defined on R
which is zero for negative values of its argument:

f (t) = 0
for any t < 0.
332
12.1.a Definition
DEFINITION 12.2 Let f (t) be a real- or complex-valued locally integrable
function defined for real values of t. The (unilateral) Laplace transform

of f ( t ) is the complex function of the complex variable p denoted fb(p) (or
F (p)) which is defined by
def
fb(p) =
f (t) e p t dt.
Whereas the function fb itself is called the Laplace tranform, the operation
f 7 fb should be called the Laplace transformation.
DEFINITION 12.3 Let f be a locally integrable function on R. The bilateral
Laplace transform of f is the function

Z +
p 7
f (t) e p t dt.
Its properties are of the same kind as those of the unilateral transform, with
variants that the reader may easily derive by herself.
The bilateral Laplace transform of H (t) f (t) is simply the Laplace transform of the function f (t). In the remainder of this chapter, only the unilateral
transform will be considered.
Remark 12.4 If f is a function zero for t < 0 which has a Fourier transform, then the relation

fb(i) = fe
with R
2
holds. The Laplace transform is an extension of the notion of Fourier transform. More
precisely, fb(x + i) is, up to the change of variable = 2, the Fourier transform of t 7
f (t) e x t evaluated at /2.
Sometimes the following definition is found in the literature:

DEFINITION 12.5 (Original, image) An original is any function of a real
variable which is locally integrable, with real or complex values, such that
i) f is causal ( f (t) = 0 for all t < 0);

ii) f (t) does not increase faster than any exponential function, that is,
there exist constants M > 0 and s R such that

f (t) M e st
for all t R.
The Laplace transform of an original is the image of the original.
333
Pierre Simon, marquis de Laplace (17491827), befriended

dAlembert (p. 415) at eighteen, found a position as professor
of mathematics at the Paris military school, and then became a
professor at the cole Polytechnique. He was passionate about
astronomy and developed many concepts and useful tools for
the study of mechanics, in particular in the field of differential
equations. He studied cosmology and cosmogony (the primitive nebula of Laplace), rediscovered the forgotten works of
Bayes (p. 518) and was the first to compute the integral of the
gaussian. When Napolon asked why he did not mention God
in his Trait de mcanique cleste, his answer was: Je nai pas eu
besoin de cette hypothse (I did not need this assumption).
The following notation are used: if f (t) is an original, its Laplace transform is denoted fb(p) or F (p). The symbol is also in widespread use, used
as follows: f (t) F (p). (In certain russian books, the notation f (t) F (p)
also appears.)
12.1.b Integrability
In the remainder of this chapter, we write p = x + i.
Let f be a locally integrable function. We are looking for the domain
where its Laplace transform fb(p) is defined. Note first that integrability of
t 7 f (t) e p t is equivalent to integrability of t 7 f (t) e x t .
Moreover, if this function is integrable for some x0 , then it is also integrable for any x > x0 , since

f (t) e x t = f (t) e x0 t e (x0 x)t f (t) e x0 t .
It follows that the set of complex numbers p, where t 7 f (t) e p t is integrable, either is empty or is a (right) half-plane in the complex plane, or
indeed is the entire complex plane C.
DEFINITION 12.6 The convergence abscissa of the original function f is the
lower bound of all real numbers x for which the function above is integrable:
n
o

def
= inf x R ; t 7 f (t)e x t is integrable .
From the previous reasoning, we deduce the following result:
PROPOSITION 12.7 Let f (t) be an original and let be its convergence abscissa.
Denote as before p = x + i. Then
i) for x < , the function t 7 f (t) e p t , is not integrable;

ii) for x > , this function is integrable;
334
iii) for x = , it may, or may not, be Lebesgue-integrable. If it is not Lebesgueintegrable, the improper integral
Z +
ZR
p
t
f (t) e
dt = lim
f (t) e p t dt
R+ R
R
may still converge for certain values of p = +i, which provides an extension
of the Laplace transform of f to those values of p.
Example 12.8 Consider the Heaviside function. Its Laplace transform is given by
b (p) =
H
H (t) e p t dt =
1 p t +
e
,
0
p
so that the convergence abscissa of H is equal to 0, and for all p C such that Re(p) > 0, we
b (p) = 1/p. The domain where H
b is defined is thus an open half-plane.
have H
b
Notice that, if the function H is extended to the imaginary axis (except the origin) by
defining, for any nonzero , H (i) = 1/i, we obtain a function which is close to the
Fourier transform of H , the latter being equal to
e () = pv 1 + ,
H
2i
2
which is the same as pv(1/2i) on R . This resemblance will be explained in Section 12.4.e,
page 345.
Example 12.9 Let f (t) = 1/(1 + t 2 ). Then we have
F (p) =
ept
dt,
1 + t2
which is defined for all p C such that Re(p) 0. The domain of definition of F is therefore
a closed half-plane.
Example 12.10 Let a C be an arbitrary complex number, and consider the Laplace transform
of the function t 7 H (t) e at . We have
Z +
Z +
H (t) e at
e at e p t dt =
e ( pa)t dt,
0
so that the convergence abscissa is = Re(a), and for all p C such that Re(p) > Re(a), the
Laplace transform of H (t) e at is 1/(p a).
THEOREM 12.11 In the open half-plane on the right of the abcissa of convergence,
the Laplace transform fb is holomorphic and hence also analytic.
Proof. Let x R with x > . Then, for any m N, the function t 7 t m f (t)e p t is
integrable: since x > , the real number y = (x + )/2 satisfies < y < x; moreover,
t 7 f (t) e yt is integrable, and we may write

t m f (t) e x t = f (t)e yt t m e (xa)t/2 ,
where the first factor is integrable while the second is bounded. The derivative of F
can be computed using the theorem of differentiation under the integral sign, and is
given by
Z
dF
= F (p) =
(t) f (t) e p t dt.
dp
0

i
335
i
<0
x
(a)
(b)
=0
>0
x
(c)
(d)
Fig. 12.1 Convergence abcissas for unilateral Laplace transforms. (a) Function with
bounded support on the right. (b) Function rapidly decaying on the right.
(c) Tempered function. (d) Rapidly increasing function. In the non-gray open
set, the Laplace transform is holomorphic.
Since F is differentiable in the complex sense at any point of the open half-plane
{p C ; Re p > }, it is an analytic function.
Notice that the convergence abcissa for t 7 t f (t) is the same as that for f . Also,
an obvious induction shows that
Z
F (n) (p) =
(t n ) f (t) e p t dt.
0
We will admit the following result, which is summarized in Figure 12.1:

PROPOSITION 12.12 Let f be an original.
i) If the function f (t) has bounded support on the right, then = .
ii) If f (t) decays rapidly on the right, then < 0 and fb(p) is holomorphic in a half-plane containing the imaginary axis (in particular, the Fourier
transform of f exists).
iii) If f (t) is tempered (in the sense of distributions) and does not tend to 0 at
infinity, then = 0; its Fourier transform does not necessarily exist in the sense
of functions, but it exists in the sens of distributions.
iv) Finally, if f (t) increases rapidly at infinity, then 0 < +; then fb (if
it exists) is holomorphic on a half-plane not containing the imaginary axis and
the Fourier transform of f does not exist.
Example 12.13 The convergence abcissa for the rectangle function is and we have
b
(p)
=
1/2
e p t dt =
1 e p/2
,
p
336
b is bounded in a
which is analytic on C; the singularity at 0 is indeed an artificial one, since
b has the power series expansion
neighborhood of 0;

X
(1)n+1
b
pn.
(p)
=
(n + 1)! 2n+1
n=0
Any nonzero polynomial function or rational function has convergence abcissa equal to 0.
2
Exercise 12.1 Show that the function g : t 7 e t does not have a Laplace transform (that
is, = +).
12.1.c Properties of the Laplace transform

PROPOSITION 12.14 Let f (t) be an original and F (p) its image by the Laplace
transform. For x > , the function F (x + i) tends to 0 as .
Proof. This is an easy consequence of the Riemann-Lebesgue lemma (see page 282):
Z +

lim
f (t) e x t e i t dt = 0.
We know that, on the right of the abscissa of convergence, the Laplace

transform is holomorphic. What happens elsewhere?
Take, for instance, f (t) = e 3t . The convergence abscissa is = 3. Its
Laplace transform is F (p) = 1/(p 3) which is holomorphic everywhere
except at p = 3. More precisely, F is defined in the convergence half-plane
{z C ; Re(z) > 3}, but it can be analytically continued to the plane with
the single point {3} removed, namely, to C \ {3}. It often happens that the
Laplace transform of an original is a meromorphic function.1 Thus we will show
that
p
;
cos t 2
p +1
this image function is meromorphic and has poles at p = i. It is represented
in Figure 12.2, together with its analytic continuation to C \ {i, i}.
12.2
Inversion
It is possible to find an inversion formula for the Laplace transform using
our knowledge of the Fourier inversion formula. Indeed, if F (p) is the Laplace
transform of the function f (t), we have, for any x > ,
Z +
F (x + 2i) =
H (t) f (t) e x t e 2it dt,
1
This is not always the case. It may have branch points, as in the case of the Laplace
transform of t 7 H (t)/(t + 1).
Inversion
1
y
337
0
1
1
0
0
x
Fig. 12.2 Representation of the Laplace transform of the function x 7 cos x. The zaxis is the modulus of the function p 7 p/(p 2 + 1). This function may be
analytically continued to the complex plane without i and i.
which shows that, for fixed x, the function 7 F (x + 2i) is the Fourier
transform of H (t) f (t) e x t . The Fourier inversion formula leads, when t is a
point where H (t) f (t) is continuous, to
Z +
x
t
H (t) e
f (t) =
F (x + 2i) e 2it d
or
H (t) f (t) =
F (x + 2i) e x t e 2it d =
1
2i
F (p) e p t d p,
Dx
where we have denoted p = x + 2i and D x is the vertical line

def
D x = {x + i ; R}
(Bromwich contour).
Remark 12.15 Note that the resulting formula gives the original as a function of the image,
for an arbitrary x > .

The result of the integral is indeed independent of the integration contour chosen since, F
being analytic on the convergence half-plane, Cauchys theorem guarantees that the difference
between the integral on D x and on D x vanishes (using the fact that lim F (p) = 0, as
given by Proposition 12.14).
The integration is performed by cutting the line D x and by closing it, on the left, by a
circle of arbitrarily large radius, containing all the singularities of F . (Sometimes it is necessary
to use a more subtle contour; see, for instance [31, 7.5].)
338
Let f be an original
and F its Laplace transform. Let be the convergence abcissa. For any point where f
is continuous, we have
Z x0 +i
1
F (p) e p t d p
for any x0 > .
f (t) =
2i x0 i
THEOREM 12.16 (Inversion of the Laplace transform)
Remark 12.17 If we start with an arbitrary function, this formula leads us back to H (t) f (t).
Can you show this using the fact that F is analytic in the half-plane x > a and decays on
vertical lines?
12.3
Elementary properties and examples
of Laplace transforms
Just as in the case of the Fourier transform, there are relations linking
the Laplace transform and operations like translation, convolution, differentiation, and integration.
12.3.a
Translation
There is a translation theorem similar to what holds for the Fourier transform,
but which requires some care here: since we are looking at the unilateral
Laplace transform, we will write explicitly the function H (t) f (t) instead of
f (t). The, writing p = x + i as usual and denoting by the convergence
abcissa for f , we have
if
then
H (t) f (t) F (p)

H (t) f (t) e a F (p + a)
for < x,
for Re(a) < x.
if
then
H (t) f (t) F (p)

H (t ) f (t ) F (p) e p
Similarly
for < x,
for < x.
Be careful to consider H (t ) f (t ), that is, the function which vanishes

for t < . There is, a priori, no relation between the Laplace transform of the
function H (t) f (t ) and that of H (t) f (t).
Elementary properties and examples of Laplace transforms
339
12.3.b Convolution
For two causal functions f and g, we define the convolution product, if it exists,
in the same manner as was done in 7.64 on page 211, and the definition
simplifies here to
Z +
Zt
f g(t) =
f (s) g(t s) ds =
f (s) g(t s) ds.
0
In particular, the quantity f g(t) depends only on the values of f and g

between the times 0 and t. This convolution, although formally identical to
what occurred with Fourier transforms, still has a causal aspect. To remember
this causality, we sometimes write explicitly H f H g instead of f g.
As for the Fourier transform, there is a convolution theorem, the proof of
which is left as an exercise:
THEOREM 12.18 Let f and g be originals with convergence abcissas respectively
equal to and .
If
f (t) F (p)
and
g(t) G (p)
then [H f H g](t) F (p) G (p)
for < x,
for < x,
for max(, ) < x.
Similarly, we have:
if
and
then
f (t) F (p)
g(t) G (p)
Z x0 +i
1
f (t) g(t)
F (s) G (p s) ds
2i x0 i
for < x,
for < x,

< x0 ,
for
+ < x0 .
12.3.c Differentiation and integration

Let us now try to find the Laplace transform of the derivative of a function
in the sense of functions. Take therefore an original function f and denote by F
its Laplace transform:
H (t) f (t) F (p).
We assume that f is differentiable in the sense of functions; then we have
Z +
Z +
h
i+
H (t) f (t)
f (t) e p t dt =
p f (t) e p t dt + f (t) e p t
.
0
One shows (see Exercise 12.6) that lim f (t) e p t = 0 so that the boundary
t+
term is equal to f (0) or, more precisely, to lim+ f (t), which is also denoted
f (0+ ).
t0
340
THEOREM 12.19 (Laplace transform and differentiation) Let f be an original
which is continuous, differentiable, and such that its derivative is also an original.
Then, if f (t) F (p), the Laplace transform of the derivative of f in the sense of
functions is given by
d
H (t)
f (t) p F (p) f (0+ ).
dt
Similarly, if f can be differentiated n times and if its n-th derivative if an original,
we have
dn
H (t) n f (t) p n F (p) p n1 f (0+ ) p n2 f (0+ ) f (n1) (0+ ).
dt
Example 12.20 Let f (t) = 1 + e t , which is an original with convergence abscissa = 0, and
with Laplace transform equal to
1
1
+
.
p
p +1
The function f is differentiable, with derivative g : t 7 e t , which is an original with
convergence abscissa = 1. Then we have
p
1
p F (p) f (0+ ) = 1 +
2=
= G (p).
p +1
p +1
1 + e t
It is very important to remark that the derivative is considered in the sense

of functions and the derivative at 0 must not involve a Dirac distribution .
This is why the term H (t) was explicitly included outside the derivative sign.
We will see later what happens if the derivative is considered in the sense of
distributions.
There is also a converse result:
THEOREM 12.21 Let f (t) be an original with convergence abscissa . Then t n f (t)
is an original with the same convergence abscissa. If f (t) F (p), then we have
dn
(t)n f (t) n F (p).
dp
Similarly, there is the following integration result:
THEOREM 12.22 (Laplace transform and integration)
Let f be an original,
with convergence abscissa equal to .

If
then
t
0
f (t) F (p)
F (p)
f (u) du
p
for < x,
for max(, 0) < x.
Remark 12.23 This may be deduced from the convolution theorem since

f (u) du = f H (t).
Elementary properties and examples of Laplace transforms
341
Conversely, we have
THEOREM 12.24 Let f be an original.
If
then
f (t) F (p)
Z
f (t)
F (z) dz
t
p
for < x,
for sup(, 0) < x.
In this last relation, the path of integration joins the point p at infinity, in
any direction in the half-plane of convergence. The result of the integration
is independent of the chosen path since F is holomorphic in this half-plane.
It is still necessary that F (p) decay faster than 1/ p at infinity (this is true
for nice functions, but will not hold in general for the Laplace transform of
distributions).
12.3.d Examples
We have already shown that
H (t)
1
p
for 0 < Re(p).
Using the integration and translation theorems, we obtain

1
pa
n1
t
1
H (t) e a t
(n 1)!
(p a)n
H (t) e a t
for Re(a) < Re(p),

for Re(a) < Re(p).
Using a simple fraction expansion, this provides a way to find the original
for the Laplace transform of any rational function. Note (although we do not
prove it) that the formula remains valid if we take a value of n which is not
an integer; then the factor (n 1)! in the formula must be replaced by (n),
where is the Euler function (see page 154).
In particular, we have:
H (t) e it
1
p i
and
H (t) e it
1
,
p + i
By linearity of the Laplace transform, we deduce that

p
H (t) cos(t) 2
H (t) sin(t) 2
,
2
p +
p + 2
0 < Re(p).
0 < Re(p).
Similarly,
H (t) cosh(t)
p
p 2 2
H (t) sinh(t)
p 2 2
0 < Re(p).
342
Remark 12.25 Be careful not to use the linearity to state something like:
H (t) cos t = Re H (t) e
i t
Re
1
p i
since taking the real part of some expression is certainly not a C-linear operation.
12.4
Laplace transform of distributions
12.4.a Definition
Since we are only discussing the unilateral Laplace transform (or, what
amounts to the same thing, the Laplace transform of causal functions, zero
for t < 0), we will only generalize the Laplace transformation to distributions
with support bounded on the left, that is, distributions in the space D+ .
Recall that there is a convolution algebra structure defined on this space
namely, the product T S exists for all T , S D+ .
DEFINITION 12.26 Let T be a distribution with support bounded on the left.
If there exists an R such that, for all x > , the distribution e x t T (t) is
tempered, then the Laplace transform of T is defined by the formula

def
T (t) Tb (p) = T (t), e p t .
Tb is a function. The lower bound of the possible values of is called the

convergence abscissa.
The fundamental example is the following:
PROPOSITION 12.27 (Laplace transform of ) The Dirac distribution being tem-
pered and with support bounded on the left, the Laplace transform of exists and is
equal to the constant function 1, with convergence abscissa = .
Proof. We have e t = for
all R, and is tempered. This proves that

= ; furthermore, we have (t), e p t = e 0 = 1.
12.4.b Properties
Let T be a distribution with support bounded on the left with Laplace transform Tb. Then the function Tb is holomorphic at any point p C such that
Re(p) > .
In addition, the following result may be shown:
343
PROPOSITION 12.28 (Sufficent condition for a Laplace transform) Any function F (p) which is holomorphic and bounded in modulus by a polynomial function
of p in a half-plane Re(p) > is the Laplace transform of a distribution in D+ .
The following important result then holds:

THEOREM 12.29 (Laplace transform and convolution) In the convolution alge-
bra D+ , the fundamental relation between convolution of distributions and product of

Laplace transforms holds:
b (p).
T U Tb (p) U
We may then state a differentiation theorem for distributions:

be a
THEOREM 12.30 (Laplace transform and differentiation) Let T D+
distribution with Laplace transform F (p), and let T denote its derivative in the sense
of distributions. Then T has a Laplace transform given by
T (t) p Tb (p).
Remark 12.31 This result is compatible with Theorem 12.19: let f be a continuous function
on R (hence locally integrable); the function t 7 H (t) f (t) is a priori discontinuous

at 0.

The derivative of H (t) f (t) is, in the sense of distributions, equal to H (t) f (t) + f (0) (t).
Applying the differentiation theorem 12.19 for functions says that

H (t) f (t) p F (p) f (0).
Since in addition we have f (0) (t) f (0), it follows that

H (t) f (t) p F (p),
which is the statement in the differentiation theorem 12.30 in the sense of distributions.
Also, the following results remain valid for the Laplace transform of distributions:
have Laplace transform T
b (p) with convergence
THEOREM 12.32 Let T (t) D+
abscissa . For , a R, we have
T (t ) Tb (p) e p
T (t) e a Tb (p + a)
HT
Tb (p)
p
for < x,
for Re(a) < x,

for max(, 0) < x.
Remark 12.33 The last equation is the translation of the integration theorem, since H f is a
primitive of the distribution f .
344
12.4.c Examples
Using the differentiation and translation theorems for the Laplace transform,
we find:
THEOREM 12.34 The Dirac distribution and its derivatives have convergence
abscissas = and
(t) 1,
(t) p,
(n) (t) p n ,
(t ) e p .
Putting
def
X+ (t) =
n=0
(t n),
it follows by linearity (and continuity) that

1
X+ (t)
.
1 ep
The convergence abscissa is = 0.
It is possible to write symbolically X+ (t) = H (t) X(t), although in
principle this is an ambiguous notation: indeed, H is not continuous at 0,
so H (t) (t) is not well defined. By abuse of notation, we will consider that
H (t) (t) = (t).
cos(t) X+ (t)
1
1 ep
and
e t X+ (t)
1
.
1 e 1 p
12.4.d The z-transform

Let f be a function of C class. We then have
X
X
H (t) f (t) X(t) =
f (n) (t n)
f (n) e n p .
n=0
ep,
Putting
z =
P
f
(n)
zn.
n=0
n=0
the Laplace transform thus obtained may also be written
DEFINITION 12.35 The z-transform of a function t 7 f (t) is the power
series
z 7
f (n) z n .
n=0
The z-transform of a function is used in signal theory, when it is necessary to

work with sampling (we do not know f (t) for all real t, but only for integer
values).
345
12.4.e Relation between Laplace and Fourier transforms

Consider a distribution f D+ , with Laplace transform F (p) and convergence abscissa .
If > 0, then f (t) is not a tempered distribution; it does not a priori
have a well-defined Fourier transform.
If < 0, then f (t) has a Fourier transform; it is linked to the Laplace
transform by the relation
fe() = F (2i).
If = 0, which is the most interesting case, it is possible to show that

F (p) is a meromorphic function which has poles only on the imaginary axis
(remember that although F is not defined on the left half-plane of the complex
plane, it may possibly have an analytic continuation to a larger domain). It
is not possible to consider directly the function F (2i), which is not locally
integrable and consequently does not correctly define a distribution the
problems occur of course at the poles. Now expand F (p) as the sum of a
holomorphic function and rational functions with poles on the imaginary
axis:
X
n
F (p) = G (p) +
,
(p in )mn
nI
with G holomorphic for Re(p) 0. The integers mn are the orders of the
poles in . By means of the inverse Laplace transform, we derive
f (t) = g(t) + H (t)
n e in t
nI
t mn 1
.
(mn 1)!
The Fourier transform of this function may be computed, as long as g has

a Fourier transform. Denoting n = 2in , and using the finite part
distribution (see Exercises 8.6 on page 241 and 11.5 on page 324), we obtain:
(mn 1) ( )
X
n
n
n
m ,
fe() = e
g () +
+ fp
2(2i)mn 1 (mn 1)!
2i( n ) n
nI
that is,
fe() = fp F (2i) +
or also
n
m
2(2i) n 1 (mn
nI
1)!
(mn 1) ( n ),

X (i)mn 1 n
fe
= fp F (i) +
(mn 1) ( n ),
2
2(m
1)!
n
nI
(12.1)
where fp is the finite part (for multiple poles), or the Cauchy principal value
(for simple poles).
346
Exercise 12.3 Compute this way the Fourier transform of the Heaviside distribution and of
H (x) x n (see Exercise 11.5 on page 324).
12.5
Physical applications, the Cauchy problem
The Laplace transform may be used to solve problems linked to the evolution of a system, when the initial condition is known. Suppose we have a
linear system of differential equations describing the evolution of a physical
system for which the equations of motion are supposed to be known. We
are looking for the solutions t 7 f (t), where the value of the function and
of a certain number of its derivatives at t = 0 are imposed. The Cauchy
problem is the data consisting of the differential equation (or the system of
differential equations) together with the initial conditions.
12.5.a Importance of the Cauchy problem

The physically essential notion which is also philosophically fascinating
of determinism is founded on the more mathematical concept of the Cauchy
problem. A system is deterministic if the knowledge, at a given time t = 0, of
the various parameters characterizing this system (typically, the positions and
velocities of the particles that it is made of, in the case of a classical system),
determines theoretically the future (or past) evolution of the system without
ambiguity.
In other words, the differential equations that describe the dynamics of the
system must have a solution, for a given initial condition, and this solution
must be unique.
A whole branch of mathematical physics deals with the proof of existence
and uniqueness theorems for solutions of physical systems (often simplified).
But it should not be forgotten, taking as excuse that this kind of work is
highly technical (and thus often arid-looking), that this is the key to the
precise notion of determinism.2
2
Many quarrels between scientists have been based on this concept. It is to be seen already
in Newtons equations and appears, in particular, later with Laplaces intuition of the deterministic character of classical mechanics (the famous I did not need this assumption). These
disagreements continued into the last century: Einstein, with his God does not play dice, and
the ever-lively discussions between partisans of various interpretations of quantum mechanics.
For a theatrical approach to this subject, the reader is invited to go see (or to read) the play
Arcadia by Tom Stoppard [88].
347
12.5.b A simple example

The Laplace transform can be used to solve a differential equation, taking
initial conditions into account, in a more natural manner than with the
Fourier transform.
The best is to start with a simple example. Suppose we want to solve the
following Cauchy problem:

x + x = 2 cos t,
(C )
x(0) = 0 x(0) = 1.
We assume that the function x has a Laplace transform: x(t) X (p). Then
the second derivative of x has the following Laplace transform, which incorporates the initial conditions:
x(t) p 2 X (p) p x(0) x(0) = p 2 X (p) + 1,
while
p
.
+1
The differential equation (with the given initial conditions) can therefore be
written simply in the form
cos t
p2
p 2 X (p) + 1 + X (p) =
2p
,
+1
p2
which is solved, in the sense of functions3, by putting

X (p) =
2p
1
2
.
(p 2 + 1)2
p +1
There only remains to find the original of X (p). The second term on the
right-hand side of the preceeding relation is the image of the sine function:
sin t 1/(p 2 + 1). Moreover, using Theorem 12.21, we obtain

d
1
2p
t sin t
= 2
,
2
dp p + 1
(p + 1)2
which gives the solution X (p) t sin t sin t = (t 1) sin t.
It is easy to check that x(t) = (t 1) sin t is indeed a solution of the stated
Cauchy problem.
The reader is referred to the table on page 614 for a list of the principal
Laplace transforms.
3
One can see the advantage of working with Laplace transforms, which are functions,
instead of Fourier transforms, which are distributions.
348
12.5.c Dynamics of the electromagnetic field without sources

The example treated here is somewhat more involved.
Consider the free evolution of the electromagnetic field in the space R3 , in
the absence of any source. For reasons of simplicity, it is easier to study the
evolution of the potential four-vector A( r, t) in Lorenz gauge instead of the
evolution of the electric and magnetic fields. The initial conditions at t = 0
are given by the data of the field A( r, 0) = A0 ( r ) and of its first derivative
0 ( r ).
with respect to time, which we denote by A
The equations which determine the field are therefore
for all r R3 and t > 0,

A( r, t) = 0
A( r , 0) = A0 ( r ) for all r R3 ,
( r , 0) = A
0 ( r ) for all r R3 .
A
e ( k, t) its
We denote by A ( r, p) the Laplace transform of A( r, t), and by A
Fourier transform with respect to the space variable. We also denote by
f( k, p) the double transform (Laplace transform for the variable t and
A
Fourier transform for the variable r ).
We have, by the differentiation theorem 12.19 :
A( r , t) A ( r , p),
A( r , t) pA
A ( r, p) A( r , 0),
t
2
( r , 0),
A( r , t) p 2A ( r , p) p A( r , 0) A
t2
and the evolution of the potential is therefore described, in the Laplace world,
by the equation

2
p
1
p
A ( r , p) = 2 A0 ( r ) + A
( r).
2
c
c
c 0
Take then the Fourier transform for the space variable, which gives
2

p
pe
1e
2 f
+ k A ( k, p) = 2 A
0 ( k) + A0 ( k),
2
c
c
c
p
1
e
0 ( k),
e ( k) +
A
A
2 2 0
2
+k c
p + k2 c 2
from which we can now take the inverse Laplace transform to derive, for any
t 0, that
e
e ( k, t) = A
e 0 ( k) H (t) cos(k c t) + A
0 ( k) H (t) sin(k c t) .
A
kc
Finally we need only compute the inverse Fourier of t 7 (sin k c t)/k c and
t 7 cos(k c t).
or
f( k, p) =
A
p2
349
THEOREM 12.36 The inverse Fourier transforms of those functions are given by
F 1
and
Z

sin(k c t)
1
sin(k c t) i k r 3
d k
=
e
3
kc
(2)
kc
i
1 h
=
(r c t) (r + c t)
4c r
1
(r c t) if t > 0,
= 4c r
0
if t = 0,
F 1 [cos(k c t)]
1
(r
4r
( r )
c t) if t > 0,
if t = 0.
Proof. See Appendix D, page 603. The derivative (v c t) is assumed to be a radial

derivative.
The expression for the evolution of the field is then obtained by this last
inverse Fourier transform using the the convolution theorem. Denoting by
B( r ; c t) the sphere centered at the point r with radius c t (i.e., the boundary
of the corresponding ball), we derive for t > 0 the formula
ZZ
ZZ
1
1
A0 ( r ) 2
( r ) d2 r .
A( r , t) =
d
r
+
A
4 B( r;c t) n | r r |
4c 2 t B( r;c t) 0
where / n is the exterior normal derivative (in the direction the vector
exterior to the surface).
Notice that the value of the field at point r and time t is expressed in
terms of the initial values of the field and its derivative at points r such that
| r r | = c t (this is what appears in the integration domain B( r ; c t)).
Note also that if we let t tend to 0 in this formula, the second term tends
to 0, whereas the first one tends to A0 ( r) if A0 is sufficiently regular (it is
enough that it be of C 1 class here).
Exercise 12.4 As an exercise, the reader may show that the previous formulas may be recov-
ered using the Kirchhoff integral presented in Problem 3 on page 244.

Start by showing that, if S is an arbitrary (smooth) surface, and if ( r , t) is a scalar field
such that = 0, then by decomposing in monochromatic waves via a Fourier transform
with respect to t, and using the results of Problem 3, we get
ZZ

1
1
1 q
1
( r , t) =
[]
+
+
d2 S ,
4 S
n q
cq n t
q n
where q is the distance from r to the point of surface which is the integration variable, and
/ n is the exterior normal derivative. The quantities between curly brackets are evaluated at
the delayed time t q/c .
To conclude, take S = B( r ; c t) and show that the last formula recovers the framed
equation above.
350
We can also express the potential evolving from t = 0 in presence of a

current four-vector j( r , t) which is assumed to be known at any time, namely,

ZZZ
j r , t 1c | r r |
A( r , t) =
d3 r
| r r |
B( r;c t)
ZZ
ZZ
1
A0 ( r ) 2
1
( r ) d2 r .
A
d r +
4 B( r;c t ) n | r r |
4c 2 t B( r;c t ) 0
In this last expression, the first term, produced by the sources, is given by a
volume integral; it is therefore clear that the current j( r , t) may be a distribution instead of a function.
Note that to be physically pertinent, we must impose that the inital conditions (the potential at t = 0 and its first derivative) must be compatible with
the initial conditions concerning the charge and current distribution.
Exercises
351
EXERCISES
Exercise 12.5 Let f (t) be an original which is periodic with period a. Show that its Laplace
transform may be written

1
F (p) =
1 e pa
f (t) e p t dt.
Exercise 12.6 We prove here an intermediate result used in the proof of Theorem 12.19.
i) Let h : R R be a differentiable function such that h and h are integrable. Show that
lim x+ h(x) = 0.
ii) Deduce that if f is an original with convergence abscissa , such that f is differentiable,
and if f is also an original with convergence abscissa , then for any complex number
p such that Re(p) > max(, ), we have
lim f (t) e p t = 0.
t+
Exercise 12.7 Let J0 be the Bessel function of order 0. Show that
t 0,
(See the table on page 614.)
J0 (s) J0 (t s) ds = sin t.
Exercise 12.8 Prove the following results, and draw the original functions in each case (k is
a given real number):
1 k p
e ,
p
a)
H (t k)
b)
(t k) H (t k)
1 k p
e ,
p2
c)
H (t) H (t k)
1 e k p
,
p
d)
n=0
H (t) + 2
e)
H (t nk)
(1)n H (t 2nk)
n=1
f)
(1)n H (t nk)
n=0
g)
t H (t) + 2
(1)n (t 2nk) H (t 2nk)
n=1
1
,
p(1 e k p )
tanh(k p)
,
p
1
p(1 + e k p )
1
tanh(k p).
p2
These results are particularly useful in electronic and signal theory.

Exercise 12.9 Consider a mechanical system given by a dampened harmonic oscillator, de-
scribed by the equation x 2

x + x = f (t), where f (t) is an outside perturbation. We assume
that the system is initally at rest, that is, x(0) = x(0) = 0, and that the outside perturbation
has the form
5
X
f (t) =
(1)n H (t n) with > 0.
n=0
Draw the graph of f , then find the solution t 7 x(t) to the given Cauchy problem.
352
SOLUTIONS
def
Solution of exercise 12.5. Put =
f
[0,a]
. Then we have f = H (t) X(t/a). Apply
then the convolution formula, remembering that H (t) X(t) 1/(1 e p ).

i) Since f is integrable, the function
f has a limit at + (indeed, for any real number x
Rx
we have x: f (x) = f (0) + 0 f (t) dt and this last expression, by definition, has a limit
as x tends to +).
ii) By definition of the convergence abscissa, for Re(p) > , we know that t 7 f (t) e p t
is integrable. Moreover, its derivative is then t 7 f (t) e p t p f (t) e p t , which is
also integrable if Re(p) is larger than and . Apply then the previous result to the
function h : t 7 f (t) e p t .
Solution of exercise 12.7. The Laplace transform of J0 is given by

H (t) J0 (t) p
1
p2 + 1
We want to find the convolution H (t) J0 (t) with itself, which has image
1
1
1
H (t) J0(t) H (t) J0 (t) p
p
= 2
,
2
2
p +1
p +1
p +1
the original of which is H (t) sin t.
Note that the formula stated is still valid for t 0 by the parity of J0 .
Solution of exercise 12.8. The graphs are
1
k
1
k
1
2k
4k
353
The formulas are obtained by simple application of the various properties of the Laplace
transform and summations.

x(t) =
5
X
n=0

(1)n 1 + e tn (t n 1) H (t n).
Chapter
13
Physical applications of the

Fourier transform
13.1
Justification of sinusoidal regime analysis
When studying a linear system1 it is customary to reason as follows: consider
a situation where the signal is purely monochromatic, and where the response
to an arbitrary signal is the sum of responses to the various frequencies that
make up the signal. The sinusoidal analysis is then justified by stating that
If the system has for input a sinusoidal signal with given frequency ,
then the output signal is necessarily sinusoidal with the same frequency .
For definiteness, let us look at the very pedestrian example of an electrical
circuit of RLC type (Resistors, Inductors, Capacitors). Assume that the input
signal is the circuit voltage t 7 u(t), and the output signal is the resistor
voltage t 7 v(t) (i.e., up to the factor R, the intensity of current in the
circuit).
1
For instance, a mechanical system with springs and masses, or an electric circuit built out
of linear components such as resistors, capacitors or inductors, or in the study of the free
electromagnetic field, ruled by the Maxwell equations.
356
Physical applications of the Fourier transform
u(t)
v(t)
Writing down the differential equation linking u and v, we obtain the

identity (valid in the sense of distributions):
L
1
u=
+
H + v = D v,
R
RC
which may be inverted (using the whole apparatus of Fourier analysis!) to

yield an equation of the type
v(t) = T (t) u(t),
where T is the impulse response of the system, that is, T = D 1 .

Let us restrict our attention to the case where the input signal is sinusoidal,
with pulsation (or frequency related to by = 2). The output signal
is assumed to be sinusoidal, with same pulsation . The equations are much
simplified because, formally, we have
d/dt = i
(a differential operator is replaced by a simple algebraic operation). Define the
transfer function Z() as the ratio between the output signal v(t) = v0 e it
and the input signal u(t) = u0 e it :
Z() =
v(t)
v
= 0.
u(t)
u0
When the input signal is a superposition of sinusoidal signals (which is, generally speaking, always the case), the linearity of the system may be used to
find the output signal: it suffices to multiply each sinusoidal component of the
input signal by the value of the transfer function for the given pulsation and
to sum (or integrate) over the whole set of frequencies:
X
X
if u(t) =
an e in t ,
then v(t) =
an Z(n ) e in t ,
n
and if
u(t) =
a() e it d,
then
v(t) =
a() Z() e it d.
The only subtle point in this argument is to justify that the output signal
is monochromatic if the input signal is. This step is provided by the following
theorem:
THEOREM 13.1 Any exponential signal (either real t 7 e p t or complex t 7
e 2it ) is an eigenfunction of a convolution operator.
Justification of sinusoidal regime analysis
357
Indeed, assume that the input signal u and the output signal v are related
by a convolution equation v(t) = T (t) u(t), where T is either a function or
a tempered distribution. If u is of the form u(t) = e 2it , then v(t) = T (t)
e 2it is infinitely differentiable (because the exponential is) and Theorem 8.24
yields

v(t) = T (), e 2i(t) = T (), e 2i e 2it ,
which is indeed a complex exponential with the same frequency as the input signal. The
transfer function is, for the frequency , equal to Z() =

T (), e 2i or, in other words, it is the value of the Fourier transform of
the impulse response T :
Z = Te.
The quantity Z() may also be seen as the eigenvalue associated to the eigenvector
t 7 e 2it of the convolution operator T .
The same reasoning applies for a real-valued exponential. If the input
signal is of the type u(t) = e p t (where p is a real number), the output signal
can be written

v(t) = T (), e p(t) = T (), e p e p t ,

still according to Theorem 8.24. The transfer function p 7 T (), e p
is then of course the Laplace transform of the impulse response T (see Chapter 12).
Remark 13.2 (Case of a resonance) This discussion must be slightly nuanced: if Z () is infi-
nite, which means the Fourier transform of the distribution T is not regular at , one cannot
conclude that the output signal will be sinusoidal.
A simple example is given by the equation of the driven harmonic oscillator:
v + 2 v = e 2i t = u(t).
This equation can be written u = [ + 2 ] v, the solution of which, in D+ , is (see
Theorem 8.33 on page 239)
1
H (t) sin t.
The Fourier transform T (computed more easily by taking the Laplace transform evaluated at
2i, and applying the formula (12.1), page 345) is then2
1
1
1

Z () = Te() = fp 2
+
+
.
2
2
4
4i
2
4i
2
v = T u
with T (t) =
For the values = /2 and = /2, there is divergence.

2
The transfer function Z () is computed, ab initio, at the beginning of Chapter 15. This
will explain why imaginary Dirac distributions appear, in relation with the choice of a solution
of the equation u = [ + 2 ] v; for a solution in D , the Dirac spikes would have the
opposite sign. A solution which does not involve any Dirac distribution exists also, but it does
not belong to any convolution algebra.
358
In fact, this reflects the fact that the solutions of the differential equation are not sinusoidal; instead, they are of the form
1
t eit,
v(t) =
2i
at the resonance frequency. More details are given in Chapter 15.
13.2
Fourier transform of vector fields:
longitudinal and transverse fields
The Fourier transform can be extended to an n-dimensional space.
DEFINITION 13.3 The Fourier transform of a function x 7 f ( x) with
variable x in Rn is given by
fe( ) =
ZZ
Rn
f ( x) e 2i x dn x.
All the properties of the Fourier transform extend easily to this case, mutatis
mutandis; for instance, some care is required for scaling:

1
.
F f (a x) = n fe
|a|
a
One may also define, coordinate-wise, the three-dimensional Fourier transform of a vector field x 7 A( x). Using this, it is possible to define transverse
fields and longitudinal fields.
Let x 7 A( x) be a vector field, defined on R3 and with values in R3 .
This field may be decomposed into three components A = (A x , A y , Az ), each
of which is a scalar field, that is, a function from R3 to R. Denote by A x the
Fourier transform of A x , A y that of A y , and A z that of Az . Putting together
these three scalar fields into a vector field k 7 A ( k), we obtain the Fourier
transform of the original vector field:
ZZZ
A ( k) =
A( x) e 2i k x d3 x.
R3
DEFINITION 13.4 Let A( x) be a vector field x Rn with values in Rn .
Denote by k 7 A ( k) its Fourier transform. For any k Rn , decompose the

vector A ( k) in two components:
the first is parallel with k and is denoted A// ( k); it is given by the
formula
def A ( k) k
A// ( k) =
k;
k kk2
Heisenberg uncertainty relations
359
the second is orthogonal to k and is denoted A ( k); it is given by

def
A ( k) = A ( k) A// ( k).
The longitudinal component of the field A is given by the inverse Fourier
transform of the field k 7 A// ( k):

A ( k) k
def
1
A// ( x) = F
k
.
| k|2
Its transverse component is then

def
A ( x) = F 1 A ( k) = A( x) A// ( x).
DEFINITION 13.5 A vector field is transverse if its longitudinal component

is zero, and is longitudinal if its transverse component is zero.
From the differentiation theorem, we can derive the well-known result:

PROPOSITION 13.6 A transverse field satisfies the relation
A( x) = 0
at any x R3 ,
that is, its divergence is identically zero.

It is important to notice that the property of being transverse or being
longitudinal is by no means a local property of the vector field (it is not sufficient to know the value of the field at a given point, or in the neighborhood
of a point), but a property that depends on the value of the field at all points
in space (see Exercise 13.1 on page 375).
Remark 13.7 A monochromatic field, with wave vector k, is transverse if it is perpendicular
to k at any k; it is easy to have an intuitive idea of this. A general transverse field is then a
sum of monochromatic transverse fields.
13.3
In nonrelativistic quantum mechanics, the state of a spinless particle is
represented by a complex-valued function which is square integrable and
of C 2 class. More precisely, for any t R, the function (, t) : r 7 ( r , t)
is in L2 (R3 ). This function represents the probability amplitude of the particle
at the given point, that is, the probability density associated with the presence
360
of the particle at point r in space and at time t is |( r , t)|2 . The function

(, t) is normalized (in the sense of the L2 norm) to have norm 1:
ZZZ

(, t) 2 =
( r , t)2 d3 r = 1.
for any t R,
2
R3
DEFINITION 13.8 The wave function in position representation is the

square-integrable function thus defined. The vector space H of the wave
functions, with the usual scalar product, is a Hilbert space.3
DEFINITION 13.9 The position operator along the O x-axis is denoted X ,
and similarly Y and Z are the position operators along the O y- and O zdef
axes, where R = (X , Y , Z). The operator X is defined by its action on wave
functions given by
X =
with (, t) : r 7 x ( r , t),
where we write generically r = (x, y, z).
Be careful not to make the (classical!) mistake of stating that and X

are proportional, with proportionality

constant equal to x!
It is customary to denote (t) the function (, t), considered as a vector
in the vector space of wave functions.4 The scalar product of two vectors is
then given by
ZZZ

def

(t) (t) =
( r , t) ( r, t) d3 r,
(recall the chosen convention for the hermitian product in quantum mechanics: it is linear on the right and semilinear on the left), and the position operator
acts as follows:

X (t) = (t)
with ( r, t) = x ( r , t).
With each vector | is associated a linear form on the space H by means

of the scalar product, which is denoted | H . Thus, the notation |
can be interpreted in two (equivalent) ways: as the scalar product of with
(in the sense of the product in L2 ), or as the action of the linear form | on
the vector | (see Chapter 14 for a thorough discussion of this notation).
3
Only physical wave functions belong to this space. Some useful functions, such as
monochromatic waves r 7 e i k r do not; they are the so-called generalized kets, which are
discussed in Chapter
14, in particular, in Section 14.2.d.
4
In fact (t) is a vector in an abstract Hilbert space
H . However, there exists an isometry

between H and H = L2 (R3 ), that is, the vector (t) can be represented by the function
r 7 ( r , t): this is the position representation.

But, since the Fourier transform is itself
an isometry of L2 (R3 ) into itself, the vector (t) may also be represented by the function
e p, t), which is the Fourier transform of the previous function. This is called the
p 7 (
momentum representation of the wave vector.
361
DEFINITION 13.10 The average position on the x-axis of a particle charac-
terized by the wave function ( r , t) is the quantity

x (t) = (t) X (t) = (t) X (t) =
def
ZZZ

2
x ( r , t) d3 r.
In addition, to every wave function ( r , t) is associated its Fourier transform with respect to the space variable r , namely,
ZZZ
d3 r
def
( p, t) =
( r , t) e i p r/ h} p
.
(2 h})3
This function represents the probability amplitude for finding the particle
with momentum p. Of course, this function is also normalized to have norm
1, since for every t R we have
ZZZ
ZZZ

2 3

( p, t) d p =
( r , t)2 d3 r = 1
(t)(t) =
by the Parseval-Plancherel identity.
DEFINITION 13.11 The wave function in momentum representation is the
Fourier transform ( p, t) of the wave function ( r , t) of a particle.

def
DEFINITION 13.12 The momentum operator P = i h} = (P x , P y , P z ) is
the linear operator acting on differentiable functions by

Px =
or
P x | = |
with (, t) : r 7 i h}
( r, t).
x
The average momentum along the x-direction of a particle characterized by

the wave function ( r , t) is
ZZZ

def

p x (t) = (t) P x (t) = i h}
( r , t)
( r , t) d3 r .
x
From the known properties of the Fourier transform, we see that it is

possible to extend the definition of the operator P x on wave functions in
momentum representation by

def
P x = F P x F 1 [] ,
which gives explicitly, by the differentiation theorem,
Px =
with
(, t) : p 7 p x ( p, t).
362
DEFINITION 13.13 The position uncertainty of a particle characterized by
the wave function ( r , t) is the quantity x defined by the following equivalent expressions:
D
2 E
2
def
= x x2 = X | X | X 2
(x)2 = x x
ZZZ
2
2
=
x x ( r , t) d3 r = kX x k2 .
DEFINITION 13.14 The momentum uncertainty of a particle characterized
by the wave function ( r , t) is the quantity p x defined by the following

equivalent expressions:

( p x )2 = p x2 p x 2 = P | P | P 2

2
ZZZ

2
=
h} ( r , t) d3 r p x 2
x
ZZZ

2
=
p x2 ( p, t) d3 p p x 2 = kP x p x k2

Note that p x2 may also be expressed as the average value of the operator
2
P x2 = h}2 x 2 in the state |, since an easy integration by parts reveals that
ZZZ
2 def
2
2
2

P x = P x = h}
( r , t)
( r, t) d3 r
x2

2
ZZZ

2
=
h} ( r, t) d3 r.
x
To perform this integration by parts, it is necessary that the derivative of

decay fast enough, which is the case when p x exists.
Remark 13.15 The average values and the uncertainties should be compared with the expectation value and standard deviation of a random variable (see Chapter 20).
We are going to show now that there exists a relation between x and
p x . To start with, we have the following lemma:
LEMMA 13.16 Let be a wave function.
i) If the wave function is translated by the vector a R3 by putting

def
a ( r , t) = ( r a, t),
then neither x nor p x is changed.
ii) Similarly, if a translation by the quantity k R3 is performed in momentum

space, or, what amounts to the same thing, if the wave function in position
representation is multiplied by a phase factor by defining
363
def
k ( r , t) = ( r , t) e i k r / h},
then neither x nor p x is changed.
r, t) which is centered
iii) Consequently, it is possible to define a wave function (
in position and momentum (i.e., such that the average position and the average
momentum are zero) and has the same uncertainties as the original function
( r , t) with respect to position and momentum.
Proof. The average position of the translated wave function a is
ZZZ
ZZZ

2

2
x a =
x a ( r, t) d3 r =
(x a x ) ( r , t) d3 r = x a x .
The position uncertainty is then

ZZZ

2
2
2
2
2
(x) a = x a x a =
x 2 ( r a, t) d3 r x a x
ZZZ

2

=
(x a x )2 ( r, t) d3 r x2 2a x x + a 2x

= x 2 2a x x + a 2x x2 2a x x + a 2x

= x 2 x2 = (x)2 .
The momentum uncertainty of | is of course unchanged.
The proof is identical for the translates in momentum space.
Finally, the centered wave function is, for a given t, equal to

def
r, t) =
(
r r , t e i p r / h} .
THEOREM 13.17 (Heisenberg uncertainty relations) Let be a wave function

for which the average position and average momentum, as well as the position and
momentum uncertainties, are defined. Then we have
h}
x p x
(13.1)
2
(and similarly for y and z).
Moreover, equality holds only if is a gaussian function.
Proof. Translating the wave function as in Lemma 13.16, we may assume that x = 0
and p x = 0. Moreover, to simplify notation, we omit the dependency on t. We have

ZZZ
ZZZ

2
2 3

d r.
(x)2 =
x 2 ( r) d3 r
and
(p x )2 = h}2
(
r
)
x

The Cauchy-Schwarz inequality for functions in L2 states that

ZZZ
2 ZZZ
ZZZ

3

f ( r )2 d3 r
g( r )2 d3 r .
f
(
r
)
g(
r
)
d
r
Applying this inequality to f ( r ) = x ( r ) and g( r ) = x ( r ), we obtain

ZZZ
2 ZZZ

ZZZ

2 3
2 3
3
2

dr

dr
x
(
r
)
(
r
)
d
r
x
(
r)
(
r)

x

x
1
2 (x)2 (p x )2 .
h}
()
364
But, using an integration by parts, since x is real and the boundary terms cancel (see
Remark 13.19 concerning this point), we can write
ZZZ
ZZZ
ZZZ

( r )2 d3 r
x ( r )
( r) d3 r =
x ( r )
( r ) d3 r,
x
x
or
2 Re
ZZZ
x ( r )
( r) d3 r
x
ZZZ

( r )2 d3 r = 1.
Since |z| |Re(z)| for all z C, the left-hand side of () satisfies

ZZZ
2 ZZZ
2

1
3
3

Re
(
r
)
d
r
(
r
)
d
r
x
(
r
)
x
(
r)

= 4.
x
x
This immediately implies (13.1).

The Cauchy-Schwarz inequality is an equality only if the two functions considered
are proportional; but it is easy to see that X and / x can only be proportional if
is a gaussian as a function of x (this is a simple differential equation to solve).
If the usual definition of the Fourier transform is kept (with the factor 2
and without h}), the following equivalent theorem is obtained:
THEOREM 13.18 (Uncertainty relation) Let f L2 (R) be a function such that
R
x 2 | f (x)|2 dx and | f (x)|2 dx both exist. Define
Z

2
2 def 1
x =
x 2 f (x) dx,
2
k f k2
Z
Z

2
2
2 def 1
42
2
e

f (x) dx.
=
f () d =
2
2
k f k2
k f k2

1
Then we have x 2 2 16
2 with equality if and only if f is a gaussian.
the second moment
Remark 13.19 A point that was left slightly fuzzy above requires explanation: the cancellation
of the boundary terms during the integration by parts. Denote by H the space of square
integrable functions , such that X and P x are also square integrable. This space, with the
norm N defined by
N ()2 = kk22 + kX k22 + kP k22

ZZZ
ZZZ
ZZZ

2

2
( r)2 d3 r +
=
x 2 ( r) d3 r +
p x2 ( p) d3 p,
is a Hilbert space.
If H is of C class with bounded support with respect to the variable x, the proof
above (and in particular the integration by parts with the boundary terms vanishing) is correct.
A classical result of analysis (proved using convolution with a Dirac sequence) is that the
space of C functions with bounded support is dense in H . By continuity, the preceding
result then holds for all H .
Remark 13.20 Another approach of the uncertainty relations is proposed in Exercise 14.6 on
page 404.
Analytic signals
365
13.4
Analytic signals
It is customary, in electricity or optics, for instance, to work with complex
sinusoidal signals of the type e it , for which only the real part has physical
meaning. This step from a real signal f(R) (t) = cos(t) to a complex signal
f (t) = e it may be generalized in the case of a nonmonochromatic signal.
First of all, notice that the real signal may be recovered in two different
ways: as the real part of the associated complex signal, or as the average of the
signal and its complex conjugate.
We can generalize this. We consider a real signal f(R) (t) and search for
the associated complex signal. In the case f(R) (t) = cos(20 t) (with 0 > 0),
the spectrum of f(R) contains both frequencies 0 and 0 . However, this
information is redundant, since the intensity for the frequency 0 is the same
as that intensity for the frequency 0 , since f(R) is real-valued. The complex
signal f (t) = e 2it associated to the real signal f(R) is obtained as follows:
take the spectrum of f(R) , keep only the positive frequencies, multiply by two,
then reconstitute the signal.
In a general way, for a real-valued signal f(R) (t) for which the Fourier transform fe(R) () is defined, its spectrum is hermitian, that is, we have
fe(R) () = fe(R) ()
for all R;
thus, half the spectrum fe(R) suffices to recover f(R) .

If the spectrum fe(R) () does not have a singularity at the origin, or at least
if is sufficiently regular, we can consider the truncated spectrum
2H () fe ().
(R)
The inverse Fourier transform of this truncated spectrum will give the analytic
signal associated to the function f(R) . Since the inverse Fourier transform of
2H () is

i
F 1 2H () = F 1 [1 + sgn ] = (t) + pv
t
(see Theorem 11.16 on page 304), we deduce that the inverse Fourier transform
of 2H () fe(R) () is

Z

i
i
f(R) (t )
f(R) = f(R) (t) pv
dt .
F 1 2H () fe(R) () = + pv
t
t t
We make the following definition:
DEFINITION 13.21 An analytic signal is a function with causal Fourier trans-
form, that is, a function f such that the Fourier transform vanishes for negative frequencies.
366
DEFINITION 13.22 Let f(R) be a real signal with sufficiently regular Fourier
transform.5 The imaginary signal associated to f(R) is the opposite of its

Z
Hilbert transform:
1
f(R) (t )
def
f(I) (t) = pv
dt .
t t
DEFINITION 13.23 The analytic signal associated to the real signal f(R) is the
function
i
t
7 f (t) = f(R) (t) + i f(I) (t) = f(R) (t) pv
def
the spectrum of which is fe() = 2H () fe(R) ().
f (t )
dt ,
t t
The analytic signal associated with f(R) is analytic in the sense of Definition 13.21: its spectrum is identically zero for frequencies < 0. Since fe()
is causal, it has a Laplace transform F (p) defined at least on the right-hand
half-plane of the complex plane. Since the signal f is the inverse Fourier
transform of fe, it is given by
f (x) = F (2i x).
Hence there is an analytic continuation of f to the upper half-plane,6 which can

be obtained using the Laplace transform of the spectrum:
f (z) = F (2iz).
This analytic continuation is holomorphic in this upper half-plane, and therefore satisfies the Kramers-Kronig relations which are neither more nor less
than the relation defining f(I) : weve made a complete loop!
The fact that f may be continued to an analytic function on the upper half-place
justifies also the terminology analytic signal.
All this may be summarized as follows:
THEOREM 13.24 Let f : R C be a function which has a Fourier transform. If
the spectrum of f is causal, then f may be analytically continued to the complex upper
half-plane.
Conversely, if f may be analytically continued to this half-plane, and if the integral
of f along a half-circle in the upper half-plane tends to 0 as the radius tends to infinity,
then the spectrum of f is causal.
(See [79] for more details.)
5
In the sense that the product H () fe(R) () exists; this means that fe(R) cannot involve a
singular distribution at the origin 0.
6
The upper half-plane because of the similitude x 7 2i x: the image of the upper
half-plane is the right-hand half-plane.
Analytic signals
367
An example of the use of analytic signals is given by a light signal, characterized by its real amplitude A(R) (t). We have
Z +
A(R) (t) =
A(R) () e 2it d,
where A(R) () is the Fourier transform of A(R) (t). However, since A(R) (t) is real,
it is also possible to write
Z +

A(R) (t) =
a() cos () + 2i t d,
0
where for > 0 we denote 2A(R) (t)() = a() e i() with a and real.
The imaginary signal associated to A(R) is of course given by
Z +

A(I) (t) =
a() sin () + 2i t d,
0
and the analytic signal is
A(t) =
a() e i() e 2it d.

The quantity a() = 2A(R) () has the following interpretation: its square
W () = a 2 () is the power spectrum, or spectral density of the complex signal A (see Definition 13.26 on the next page). This complex spectral density
contains the energy carried by the real signal, and an equal amount of energy
carried by the imaginary signal. The total light intensity carried by the real
light signal is then equal to
Z
Z +
1 +
E =
W () d =
|A(R) ()|2 d,
2 0
namely, half of the total complex intensity (the other half is of course carried
by the imaginary signal).
368
13.5
Autocorrelation of a finite energy function
13.5.a Definition
Let f L2 (R) be a square integrable function. We know it has a Fourier
transform fe().
DEFINITION 13.25 A finite energy signal is any function f L2 (R). The
total energy of a finite energy signal is

Z

2
def
E = f (t) dt.
The Parseval identity states that

Z
Z

f (t)2 dt = fe()2 d,

2
so we can interpret fe() as an energy density per frequency interval d.
DEFINITION 13.26 The spectral density of a finite energy signal f is the

2
function 7 fe() .
In optics and quantum mechanics, for instance, the autocorrelation function of f is often of interest:
DEFINITION 13.27 Let f be a square integrable function. The autocorrela-
tion function of f is the function = f f given by (x) = f (x) f (x),

that is,
Z
def
(x) = f (t) f (t x) dt
for any x R.
13.5.b Properties
Since the autocorrelation function is the convolution of f with the conjugate
of its transpose (x) = f (x) f (x), it is possible to use the known properties
of the Fourier transform to derive some interesting properties. In particular,
THEOREM 13.28 The Fourier transform of the autocorrelation function of f is equal
to the spectral density of f :

2

F (t) = F f (x) f (x) = fe() .
Autocorrelation of a finite energy function
369
PROPOSITION 13.29 The autocorrelation function of any finite energy function f
is hermitian: (x) = (x) for all x R.

Moreover, if f is real-valued, then is an even, real-valued function.

In addition, since e 2i x 1, we have the inequality
Z
Z
2 2i x

2

(x) = fe() e
d fe() d = (0),

from which we deduce:
THEOREM 13.30 The modulus of the autocorrelation function of a function f is
maximal at the origin: |(x)| R (0) for any x R. (0) is real, positive, and
equal to the energy of f : (0) = | f |2 .
DEFINITION 13.31 The reduced autocorrelation function, also called the
self-coherence function of f , is given by

def
(x) =
(x)
.
(0)
This function (x) is thus always of modulus between 0 and 1.
13.5.c Intercorrelation
DEFINITION 13.32 The intercorrelation function of two square integrable
functions f and g is the function

def
f g (x) =
f (t) g(x t) dt.
From the Cauchy-Schwarz inequality, we derive

f g (x)2 f (0) g (0).
This brings us to the definition of the coherence function:

DEFINITION 13.33 The intercoherence function, or coherence function,
of f and g is the reduced intercorrelation function

def
f g (x) =
f g (x)
f (0) g (0)
The modulus of this coherence function is therefore also always between 0

and 1.
370
13.6
Finite power functions
13.6.a
Definitions
Up to now, the functions which have been considered are mostly square inte2
grable, that is, functions
R in 2L (R). In many physical situations (in optics, for
instance), the integral | f | represents an energy. Sometimes, however, the
objects of interest are not functions with finite energy but those with finite
power, such as f (t) = cos t.
DEFINITION 13.34 A function f : R C is a finite power function if the
limit
ZT
1
lim
T + 2T T
exists and is finite. This limit is then
signal f .

f (t)2 dt
the average power carried by the
The spectral density may also be redefined:

DEFINITION 13.35 Let f be a finite power signal. The spectral density, or
power spectrum, of f is the function defined by

2
Z

1 T
def

f (t) e 2it dt .
W () = lim

T + 2T T
An analogue of the Parseval formula holds for a finite power signal:
THEOREM 13.36 Let f (t) be a finite power signal and W () its spectral density.
Then we have
1
lim
T + 2T

f (t)2 dt =
W () d.
13.6.b Autocorrelation
For a finite power function, the previous definition 13.27 of the autocorrelation is not valid anymore, and is replaced by the following:
DEFINITION 13.37 The autocorrelation function of a finite power signal f
is given by
1
(x) = lim
T + 2T
def
T
T
f (t) f (t x) dt.
The Wiener-Khintchine theorem
S1 (t)
S (t)
371
S2 (t)
(S )
Fig. 13.1 Light interference experiment with Young slits.
The analogue of Theorem 13.28 is

THEOREM 13.38 The Fourier transform of the autocorrelation function of a finite
power signal is equal to the spectral density:

F (t) = W ().
Example 13.39 Let f (t) = e 2i 0 t . The autocorrelation function is (x) = e 2i 0 x and the
spectral density is W () = ( 0 ). Using the definition for a finite energy function, we
would have found instead that fe() = ( 0 ) and the spectral density would have to be
something like 2 ( 0 ). But, as explained in Chapter 7, the square of a Dirac distribution
does not make sense.
13.7
Application to optics:
the Wiener-Khintchine theorem
We consider an optical interference experiment with Young slits which are
lighted by a point-like light source (see Figure 13.1).
We will show that, in the case where the source is not quite monochromatic
but has a certain spectral width , the interference figure will get blurry at
points of the screen corresponding to differences in the optical path between
the two light rays which are too large. The threshold difference between
the optical paths such that the diffraction rays remain visible is called the
coherence length, and we will show that it is given by = c /.
Consider a nonmonochromatic source. We assume that the source S
of the figure emits a real signal t 7 S(R) (t), and that an analytic
complex

signal t 7 S (t) is associated with it. Hence, we have S(R) (t) = Re S (t) . The
Fourier transform of S(R) (t) is denoted S(R) ().
Moreover, we assume that the signal emitted has finite power and is characterized by a spectral power density W () for > 0. The signal is non-
372
monochromatic if the spectral density W () is not a Dirac distribution.

In the proposed experiment, we look for the light intensity arriving at the
point P , given by the time average of the square of the signal S P (t) arriving
at the point P :

I P = S P2 (t) ,
where we put
ZT

2 def
S (t) = lim 1
S P (t)2 dt.
P
T + 2T T
The signal reaching P is the sum of the two signals coming from the two slits.
The signal coming from one slit is thus the same as the signal coming from
the other, but with a time shift = L/c , where L is the path difference
between the two rays (and hence depends on the point P ):
S P (t) = S1 (t) + S2 (t) = S1 (t) + S1 (t ).
(The function S1 (t) is directly related to the signal S (t) emitted by the source.)
From this we derive
I P = S P S P = (S1 + S2 ) (S1 + S2 )
= S1 S1 + S2 S2 + 2Re S1 S2
= 2I + 2Re S1 S2 ,
denoting by I the average value S1 S1 , which is the light intensity of

the signal that passes by a single slit without interference. All information
concerning the interference phenomenon is given by the term 2Re S1 S2 ,

which is easy to compute:
S1 S2 = S1 (t) S1 (t ) = (),
where () is the autocorrelation function of the complex signal S1 , also

called here the autocoherence function.7 (Recall that is a function of the
observation point P .)
So finally the light intensity at the point P is given by
h

i
I P = 2I + 2 Re () = 2I 1 + Re ()
with
def
() =
7
()
()
=
.
(0)
I
In some optics books, the normalized autocorrelation function () = ()/(0) is also

called the autocoherence function.
The Wiener-Khintchine theorem
373
min
max
(0)
Fig. 13.2 The autocorrelation function for an almost monochromatic source, varying
with the parameter in the complex plane. The value at the origin (0) is
real and positive. The modulus of (t) varies slowly compared to its phase.
We may then estimate the visibility factor of the diffraction figure, using
the Rayleigh criterion [14]:
V =
Imax Imin
.
Imax + Imin
(13.2)
Here, Imax represents a local maximum of the light intensity, not at the point P ,
but in a close neighborhood of P ; that is, the observation point will be moved
slightly and hence the value of the parameter will change until an intensity
which is a local maximum is found, which is denoted Imax . Similarly, the
local minimum in the vicinity of P is found and is denoted Imin .
Estimating this parameter is only possible if the values max and min
corresponding to the maximal and minimal intensities close to a given point
are such that the difference max min is very small compared to the duration
of a wave train. If this is the case, the autocorrelation function will have
the shape in Figure 13.2, that is, its phase will vary very fast compared to its
modulus; in particular, the modulus of will remain almost constant between
max and min , whereas the phase will change by an amount equal to .8 For
= 0, we have of course = I , and then |()| will decrease with time.
Thus, by moving slightly P , varies little and || remains essentially
constant. Then Re() varies between || and + ||, and the visibility factor
8
This is what the physics says, not the mathematics! This phenomenon may be easily understood.
Suppose is made to vary so that it remains small during the duration of a wave train. Then,
taking a value of equal to half a period of the signal (which is almost monochromatic), the
signals S1 and S2 are simply shifted by : they are in opposition and () is real, negative,
and with absolute value almost equal to (0). If we take values of much larger, the signals
become more and more decorrelated because it may be that S1 comes from one wave train
whereas S2 comes
from another. A simple statistical model of a wave train gives an exponential
decay of ().
374
for the rays is then simply

V () = |()| =

()
I
But, according to Theorem 13.38, the autocorrelation function is given by the

inverse Fourier transform of the power spectrum:
Z +
() =
W () e 2i d.
0
The following theorem follows:

THEOREM 13.40 (Wiener-Khintchine) The visibility factor of the diffraction figure
in the Young interference experiment is equal to the modulus of the normalized Fourier
transform of the spectral density of the source evaluated for the value corresponding to
the time difference between the two rays:
R +

2i d
0 W () e

V = R +
.

W () d
0
So, take the case of an almost monochromatic source for instance, a

spectral lamp, which should emit a monochromatic light but, because of the
phenomenon of emission of light by wave trains, is characterized by a nonzero
spectral width. If we denote by the spectral width, we know that the
Fourier transform of the spectral energy density has typical width t = 1/;
for a time difference 1/, that is, for a path difference c /,
the diffraction pattern will become blurred. The quantity t = 1/ is the
coherence time and the length = c / is the coherence length.
Remark 13.41 A similar theorem, due to P. K. van Cittert and Frederik Zernike (winner
of the Nobel prize in 1953), links the visibility factor of a diffraction figure to the Fourier
transform of the spacial distribution of a monochromatic extended source. The interested
reader may consult [14] or look at Exercise 13.2.
Exercises
375
EXERCISES
Exercise
13.1 (Electromagnetic
field in Coulomb gauge) We consider an electromagnetic

field E( r , t), B( r , t) compatible with a charge and current distribution ( r, t), j ( r, t) ,
that is, such that the fields satisfy the differential equations (Maxwell equations):
B
,
t
1 E
curl B = 2
+ 0 j.
c t
div E = /0 ,
curl E =
div B = 0,
i) Show that the field B is transverse.
ii) Show that the longitudinal component of the electric field is given by the instantaneous
Coulomb potential associated with the charge distribution. What may be said from the
point of view of relativity theory?
iii) Is the separation between transverse and longitudinal components preserved by a
Lorentz transformation?
Exercise 13.2 (van Cittert-Zernike theorem) Consider a light source which is translation in-
variant along the O y-axis, characterized by a light intensity I (x) on the O x-axis, such that the
light passes through an interferential apparatus as described below. The distance between the
source and the slits is f and the distance to the screen is D.
y
a
We assume that the source may be modeled by a set a point-like sources, emitting wave
trains independently of each other, with the same frequency . Since the sources are not
coherent with respect to each other, we admit that the intensity observed at a point of the
interference figure is the sum of intensities coming from each source.9 Show:
THEOREM 13.42 (van Cittert-Zernike) The visibility factor of the source I (x) is equal to the modulus
of the Fourier transform of the normalized spatial intensity of the source for the spatial frequency k =
a/ f , where a is the distance between the interference slits and f is the distance between the source and
the system:

a
I (x)
def
|| = Ie
with
I (x) = R +
.
f
I (s) ds
The reader will find in [14] a generalization of this result, due to van Cittert (1934) and
Zernike (1938).
9
We saw in Exercise 1.3 on page 43 that it is in fact necessary to have a nonzero spectral
width before the sources can be considered incoherent.
376

SOLUTIONS
Solution of exercise 13.1. Applying the Fourier transform, the Maxwell equations become
k E = i e/0 ,
k B = 0,
B
B
,
t
1 E
E
ik B = 2
+ 0 ej.
c t
ik E =
The third of these shows that B is transverse.

Moreover, since k E = k E// ( k) = i e( k)/0 , we find that
k E ( k)
k
= i 2 e( k)/0 ,
k2
k
and hence, taking the inverse Fourier transform, we have
Z
( r ) [ r r ] 3
1
d r = E Coulomb ( r , t).
E // ( r , t) =
40
k r r k3
E// ( k) =
Thus, the longitudinal field propagates instantaneously through space, which contradicts special relativity. This indicates that the decomposition in transverse and longitudinal components
is not physical.
In fact, it is not preserved by a Lorentz transformation, since the Coulomb field is not.
Solution of exercise 13.2. The intensity at a point P with coordinate y on the screen is
given by the integral of intensities produced by I (x), for x R. An elementary interference
computation shows that, for small angles, we have
Z +

2a x
y
J (x) = K
I (x) 1 + cos
+
dx
f
D

Z +

2a x
y
= K I0 1 +
I (x) cos
+
dx ,
f
D
R +
with I0 = I (s) ds, and where I (x) = I (x)/I0 is the normalized intensity. Since I and I
are real-valued functions, we may write

Z +

2ia x
y
J ( y) = K I0 1 + Re
I (x) exp
+
dx
f
D

a
e 2i a y/D
.
= K I0 1 + Re Ie
f
If we decompose Ie(a/ f ) = || e i in terms of a phase and modulus, we obtain
2a y
J ( y) = K I0 1 + || cos
.
D
Now, if we observe the interference figure close

to the center and
vary y slowly, the intensity
J on the screen varies between K I0 1 || and K I0 1 + || , which shows that || is the
visibility factor of the interference figure, which is what was to be shown.
This result is useful in particular to measure the angular distance between two stars, or the
apparent diameter of a star, using two telescopes set up as interferometers.
Chapter
14
Bras, kets,
and all that sort of thing
14.1
Reminders about finite dimension
In this section, we consider a vector space E , real or complex, with finite dimension n, equipped with a scalar product (|) and with a given orthonormal
basis (e1 , . . . , en ). Some elementary properties concerning the scalar product,
orthonormal bases, and the adjoint of an endomorphism will be recalled;
these results should be compared with those holding in infinite dimensions,
which are quite different.
14.1.a Scalar product and representation theorem

THEOREM 14.1 (Representation theorem) The space E is isomorphic to its dual
space E . More precisely, for any linear form on E , there exists a unique a E
such that
x E
(x) = (a|x) .
Proof
P
Existence. Using the orthonormal basis (e1 , . . . , en ), define a = nj=1 (e i ) e i .
Pn
Pn
For any x = i=1 xi e i E , we have (x) = i=1 xi (e i ), hence (x) = (a|x).
Uniqueness. Assume a and a E satisfy
(a|x) = (a |x)
x E .
Bras, kets, and all that sort of thing
378
Then (a a |x) = 0 for any x E , and in particular (a a |a a ) = 0, which

implies that (a a ) = 0 and so a = a .
14.1.b Adjoint
DEFINITION 14.2 (Adjoint) Let u be an endomorphism of E . The adjoint
of u, if
it
exists, is the
unique v L (E ) such that, for any x, y E , we have
v(x) y = x u( y) . It is denoted u , and so by definition we have

x, y E
u (x) y = x u( y) .
Proof of uniqueness. Let v, w both be adjoints of u. Then we have

x, y E
x u( y) = v(x) y = w(x) y ,
which shows that v(x) w(x) E = {0} for any x E and hence v = w.
THEOREM 14.3 (Existence of the adjoint) In a euclidean or hermitian space E
(hence of finite dimension), any endomorphism admits an adjoint. Moreover, if u

L (E ) and if B is an orthonormal basis of E , the matrix representing u in the
basis B is the conjugate transpose of that of u:
matB (u ) = t matB (u).

Proof. For any x E , the map x : y 7 x u( y) is a linear form; according to
Theorem 14.1, there

exists a unique vector, which we denote u (x), such that x is the
map y 7 u (x) y .
It is easy to see that the map x 7 u (x) is linear. So it is the adjoint of u.
Another viewpoint, less elegant but more constructive, is the following.
Choose an orthonormal basis B. Notice first that if x, y E are represented by
matrices X and Y in this basis, we have (x| y) = t X Y . Notice also that if t X A Y =
t X BY for any X , Y M (C), then in fact A = B.
n1
Let now u, v L (E ), and define U = matB (u) and V = matB (v). Thanks to the
preceding remark,

v is an adjoint of u x, y E
v(x) y = x u( y)
X , Y Mn1 (R)
(V X )Y = t X (U Y )
t V = U .
Hence it suffices to define u as the endomorphism represented by the matrix t U in

the basis B.
COROLLARY 14.4 (Properties of the adjoint) Let u, v L (E ) and let C :
(u + v) = u + g ,
(u ) = u,
det u = det u,
IdE = Id E ,
rank(u ) = rank(u),
(u v) = v u ,
tr (u ) = tr (u),
Sp(u) Sp(u ).
Moreover, if u is invertible, then u is invertible and (u )1 = (u 1 ) .
Kets and bras
379
14.1.c Symmetric and hermitian endomorphisms

DEFINITION 14.5 (Hermitian/symmetric endomorphisms) Let u L (E )
be an endomorphism of E . Then u is called hermitian, symmetric, or selfadjoint if

x, y E
u(x) y = x u( y) ,
or equivalently if u = u. In other words, u is hermitian if and only if the
matrix representing u in an orthonormal basis satisfies A = tA (hermitian
matrix).
In the euclidean case (for a real vector space), the terminology used is
symmetric or self-adjoint, and amounts to saying that the matrix representing
u in an orthonormal basis is symmetric.
THEOREM 14.6 (Spectral theorem) Let u be a normal endomorphism of E , that
is, such that u u = u u. Then there exists an orthonormal basis E = (e1 , . . . , en )
such that the matrix representing u in E is diagonal.
Similarly, if A Mn (C) commutes with tA, then there exist a unitary matrix U
and a diagonal matrix D such that A = U D t U .
Since a self-adjoint endomorphism obviously commutes with its adjoint,

the following corollary follows:
COROLLARY 14.7 (Reduction of a self-adjoint endomorphism) Let u be a self-
adjoint endomorphism (i.e., hermitian in the complex case, or symmetric in the real
case). Then u can be diagonalized in an orthonormal basis.
Also the reader should be able to prove the following result as an exercise:
THEOREM 14.8 (Simultaneous reduction) If u are v two commuting self-adjoint
endomorphisms of E , then they can be diagonalized simultaneously in the same orthonormal basis.
14.2
Kets and bras
In the remainder of this chapter, H is a separable Hilbert space (i.e., admitting a countable Hilbert basis).
14.2.a
Kets | H
DEFINITION 14.9 (Kets) Following Dirac, vectors in the space H will be

generically denoted |. They are also called the kets of the space H.
A function with values in H will be denoted t 7 (t) .
380
An important Hilbert space is the space L2 (R). A function L2 (R)

will be denoted |. If this function is furthermore
a function of time t, the

function x 7 (x, t) will be denoted (t) .
Remark 14.10 (Spaces of quantum mechanics) In the setting of quantum mechanics, it is usual
to work at the same time with two or even three Hilbert spaces:
an abstract Hilbert space H, the vectors of which are denoted |;
the space L2 , in which the vectors are denoted (meaning a function x 7 (x)); this
space will be denoted L2 (R
R, d x) for added precision;
the space L2 again, but in which the vectors are denoted : p 7 (p); this space will
be denoted L2 (R
R, d p).
The square integrable functions correspond to Schrdingers formalism of wave mechanics, in
position representation in the case of the space L2 (R, dx), and in momentum representation
in the case of the space L2 (R, dp). The function p 7 (p) is the Fourier transform of the
function x 7 (x). We know that the Fourier transform gives an isometry between L2 (R, dx)
and L2 (R, dp).
When working at the same time with the abstract space H and the space L2 (R), the
distinction can be be made explicit, for instance, by writing a function in L2 (R) and | the
corresponding abstract vector. However, we will not make this distinction here.
In the case of a particle confined in an interval [0, a] (because of an infinite potential
sink), the space used is L2 [0, a] for the position representation and 2 for the momentum
representation; going from one to the other is then done by means of Fourier series.
Moreover, one should notice that all that is done here in the space L2 (R) may be easily
generalized to L2 (R3 ) or L2 (Rn ).
We will denote , H the scalar product in H. The norm of a vector |

is therefore

def

| 2 =
| , | H .
H
14.2.b Bras | H
DEFINITION 14.11 (Topological dual space) The (topological) dual space
of H, denoted H , is the space of continuous linear forms on H. Often,

we will simply write dual for topological dual.
DEFINITION 14.12 (Bras) Elements of H are called bras and are denoted
generically |. The result of the application of a bra | on a ket | is

denoted |, which is therefore a bracket.1
Example 14.13 Consider the Hilbert space L 2 (R). The map
: L2 (R) C
Z +
2
7
(x) e x dx
1
Dirac did not simply make a pun when inventing the terminology bra and ket: this
convenient notation became quickly indispensable for most physicists. Note that Dirac also
coined the words fermion and boson, and introduced the terminology of c -number, the function, and the useful notation h} = h/2.
Kets and bras
381
is linear. Defining g : x 7 e x , we have () = g, H for any L2 ; by the CauchySchwarz inequality, we have

() = g, kgk kk ,
L2
H
2
2
and hence the linear form is continuous and is an element of the topological dual of L2 (R).
More generally, if we choose a ket H, a linear form

: | 7 | , | H
is associated to using the scalar product. The Cauchy-Schwarz inequality

shows that this linear form is continuous (see Theorem A.51 on page 583).
It corresponds therefore to a bra, which is denoted |:

def
| = | , | H .
In other words:
To each ket | H is associated a bra | H .

The converse is also true, as proved by the Riesz representation theorem:
THEOREM 14.14 (Riesz representation theorem) The map | 7 | is a semi-
linear isomorphism of the Hilbert space H with its topological dual H . In other words,
for any continuous linear form |, there exists a unique vector | such that

| = | , | H .
| H
Hence there is a one-to-one correspondance between bras and kets:
To each bra | H one can associate a ket | H.
Because of this, it could be conceivable to work only with H and to forget

about
H ,stating that | is nothing more nor less than a notation for
| , | H . However, this result will not remain true when we consider

generalized bras.
When we define on H the operator norm

def

| =
| ,
sup
H
k|kH =1
the isomorphism above,
H H ,
| 7 | ,
becomes a semilinear isometry:

|

= | H .
In particular it is permissible to write kk without ambiguity whether

corresponds to a bra or a ket.
382
Remark 14.15 The semilinearity of the map above means that | = | for any C
and H. This is a result of the convention chosen; if we associate to a ket | the linear
map

| 7 | , | H ,
the result map H H is a linear isomorphism.
Counterexample 14.16 Consider again the Hilbert space H = L 2 (R). Let x0 R. Then the
map x0 : 7 (x0 ) is a linear form on the subspace D of continuous functions2 in H.

However, it does not take much effort to discover that it is not representable as a scalar
product with a square integrable function. Whats the problem? Certainly, it must be the case
that x0 is continuous?
Well, in fact, it is not! It is continuous on the space D with its topology (of uniform
convergence), but not on L2 with the hermitian topology (based on integrals, so defined up
to equality almost everywhere).
2
To see this, it suffices to define a sequence of functions n : x 7 e nx /2 . Then, with
x0 = 0 and = 0 (which is no loss of generality), we have
Z +
2
0,
kn k2H =
e nx dx =
n n
+
whereas

(n ) = 1 for all n N. So there does not exist a constant R such that
() kk for all H, which means that is not continuous.
H
14.2.c Generalized bras
We will denote also with the generic notation | the generalized bras, namely
linear forms which are not necessarily in H . Such a form may not be continuous, or it may not be defined on the whole of H (as seen in the case of ,
in counterexample 14.16); still it will be asked that it be defined at least on a
dense subspace of H.
One of the most important example is the following:
DEFINITION 14.17 (The bra x0 |) Let x0 R. We denote by x0 | the discontinuous linear form defined on the dense subspace of continuous functions
by
x0 | = (x0 ).
The bra x0 | is thus another name for the Dirac distribution (x x0 ).
Consider now a broader approach to the notion of generalized bra. For

this, let T be a tempered distribution. Since S is a dense subspace of L2 (R),
the dual S of S contains the dual of L2 (R), which is L2 (R) itself because
of the Riesz representation theorem.3 Hence we have the following situation:
S L2 (R) S .
2
Since functions in L2 are defined up to equality almost everywhere, x0 is not well defined
on the whole of H. We therefore restrict it to the dense subspace of continuous functions.
3
In fact, the dual of L2 (R) is isomorphic to L2 (R), but via this isomorphism one can
identify any continuous linear form on L2 (R) to a function in L2 (R).
Kets and bras
383
This is called a Gelfand triple. It is tempting to define the bra T | by

T | = T , ; but this last expression is linear in T , in the sense that
T , = T , , whereas the hermitian product is semilinear on the left.
To deal with this minor inconvenience, the following definition is therefore
used:
DEFINITION 14.18 (Generalized bra) Let T be a tempered distribution. The
generalized bra T | is defined on the Schwartz space S by its action on any
| S given by
def
T | = T , .
Remark 14.19 It is immediate to check that, if T is a regular distribution, associated to a
function S , then we have

T =
(x) (x) dx,
which is equal to the hermitian scalar product | or (|) in L2 (R).

If T is a real distribution (such as ), we do, however, have T | = (T |).
14.2.d Generalized kets

In the case where H = L2 (R) or L2 [a, b], it is possible to associate a generalized ket to some generalized bras. If | is a linear form defined by a regular
distribution
Z
+
| =
where is a locally integrable function, we define the associated generalized

ket |, namely, simply the function . This does not necessarily belong
to L2 and is therefore not necessarily an element of H. All that is asked is
that the integral be defined.
DEFINITION 14.20 (Ket | p and bra p|) Let p be a real number. The func-
tion
e i p x/ h}
2 h}
is not square integrable. It still defines a noncontinuous linear form on the
dense space L1 L2 :
p : x 7 p
p : L1 (R) L2 (R) C
Z +
7
p (x) (x) dx =
e(p),
where e is the Fourier transform of , with the conventions usual in quantum mechanics. This generalized bra is denoted p| = p |. The ket that
corresponds to the function p is denoted | p = | p
/ H.
384
Moreover, the notion of distribution can be seen as a generalization of

that of function. Hence, the generalized bra p| not only acts on (a subspace
of) L2 (R), associating to a function its Fourier transform evaluated at p, but
can be extended to tempered distributions. If |T is a ket formally associated
to a tempered distribution T S , then we have p| T = Te ( p ).
THEOREM 14.21 The family of kets |p is orthonormal, in the sense that for all
p, p R, we have p p = (p p ).
Proof. Apply the definition p| p =

p ( p ), and remember that the Fourier
transform of p in the sense of distributions is p .
The boundary between bra and ket is not clear. A distribution is a

linear form, thus technically a generalized bra. But as a generalization of
a function, it can formally be denoted as a ket. Thus we define the
|x0 associated to the generalized function (x x0 ). In particular, for any
function S , we have
| x0 = x0 | = (x0 ).
Also, the kets |x satisfy the orthonormality relation:

x| x0 = (x x0 ) = x (x0 ),
valid in the sense of tempered distributions that is, which makes sense only
when applied to a function in the Schwartz space.
14.2.e Id =
P
n
| n n |
Let (n )nN be an orthonormal Hilbert basis of H. We know that this property is characterized by the fact that
x H
x=
X
n=0
(n |x) n .
Using Diracs notation, this equation is written as follows:

| H
| =
n | |n ,
|n n | .
n=0
or also, putting the scalar n | after the vector |n

| H
| =
n=0
P
Or still, formally, the operator
n=0 |n n | is the identity operator. This is
what physicists call a closure relation.
For easy remembering, this will be stated as
Kets and bras
385
THEOREM 14.22 (Closure relation) A family (n )nN is a Hilbert basis of H if
and only if
Id =
X
n=0
|n n | .
This point of view allows us to recover easily another characterization of a

Hilbert basis, namely the Parseval identity. Indeed, for any | H, we have
X
X
2
kk = | = |
|n n | | =
| n n |
n=0
n=0

X

n | 2 ,
=
n=0
which can be compared with the formula d) in Theorem 9.21 on page 259.
14.2.f Generalized basis

We now extend the notion of basis of the space L2 (R), to include distributions, seen as generalizations of functions. Let (T ) be a family of tempered
distributions, indexed by a discrete or continuous set4 . To generalize the
idea of a Hilbert basis, the first idea would be to impose the condition
Z
Z
Id = |T T | d,
that is,
| = |T T | d,
where the notation d may represent an integral, a sum, or a mixture of

both. Still, there would be the problem of defining correctly the integral on
the right: the integrand is not a scalar but a vector in an infinite-dimensional
space.
It is easier to consider another characterization of a Hilbert basis: the
Parseval formula, which only involves scalar quantities.
DEFINITION 14.23 (Generalized basis) The family (T ) is a generalized
basis, or a total orthonormal system, if and only if we have
Z

2
2
| S
kk = | = T | d.
2
The family (T ) is a total system (or a complete system) if, for all S ,
we have

T , = 0 for any = 0 .
4
The choice of index is not without meaning. A generalized basis can be useful to diagonalize a self-adjoint operator which has, in addition to eigenvalues, a continuous spectrum.
386
As in the case of the usual Hilbert basis (see Theorem 9.21), a generalized
basis is a total system. However, the normalization of such a family (the
elements of which are not a priori in L2 ) is delicate.
Note that a generalized basis is not in general a Hilbert basis of L2 , for
instance, because any Hilbert basis is countable.
THEOREM 14.24 (The basis { x| ; x R }) Denote by x| = x| the generalized

bra associated to the distribution x . Then the family of distributions x| ; x R
is a generalized basis of L2 (R).
Proof. Indeed,
for
S
very definition of the L2 norm, we have
2 any R +
, by the
R +
2
2
kk2 = (x) dx = x| dx.
THEOREM 14.25 (The basis { p| ; p R }) Denote by p| = p | the generalized bra defined by the regular distributions p associated to the function
e ix/ h}
p : x 7 p
.
2 h}

Then the family p| ; p R is a generalized basis of L2 (R).
Proof. Let S . Then p| = (p);

e
the Parseval-Plancherel identity (Theorem 10.36), which is a translation of the fact that the Fourier transform is an isometry
of L2 (R), can be written as
Z +
Z +

2

(p)
p| 2 dp.
kk22 = kk
e 22 =
e dp =
Finally, it can
be shown that
the closure relation
is still valid for two generalized bases x| ; Rx R and p| ; p R . Indeed, if , L2 (R),
the relation | = | p p| d p holds: this is (once more) the ParsevalPlancherel identity! Similarly, we have
Z +
Z +
| =
(x) (x) dx =
x| x| dx
| x x| dx.
It will be convenient to state this as follows:
THEOREM 14.26 (Generalized closure relation)
Id =
|x x| dx
and
Id =
| p p| d p.
Remark 14.27 This shows that it is possible to write, formally, identities like
| =
|x x| dx,
that is,
() =
x () (x) dx.
Linear operators
387
This formula has a sense if evaluated at a point y; indeed, this yields the correct formula
Z +
( y) =
(x) x ( y) dy.
The great strength of Diracs notation is to lead, very intuitively, to many correct formulas. For
instance, writing
Z
+
| =
| p p| dp
amounts to saying that, for any x R, we have

Z +
Z +
1
e i p x/ h} (p)
e dp,
(x) = x| =
x| p p| dp = p
2 h}
which is none other than the inverse Fourier transform formula!
14.3
Linear operators
We will only describe here the simplest properties of linear operators acting
on Hilbert spaces. This is a very complex and subtle theory, which requires by
itself a full volume. The interested reader can read, for instance, the books [4,
95] or [71], which deals with quantum mechanics in Hilbert spaces. We will
define bounded operators, closed operators, eigenvalues and eigenvectors, and
adjoints of operators.
14.3.a
Operators
DEFINITION 14.28
(Linear operator) A linear operator on a Hilbert space H
is a pair A , D A , where
D A is a (not necessarily closed) subspace of H, called the domain of A;

A is a linear map defined on D A and with values in H.
The image of D A by A is called the image of A, and denoted Im( A).

Finally the image of an element H by A will be denoted either A() or,
more often, simply A.
Remark 14.29 To say that two operators A and B are equal means, in consequence, that:
their domains coincide: D A = D B ;

they act identically on their domains: A = B for any D A .
Very often, the action is only used as a shorthand for the whole operator; for instance, the
position operator in quantum mechanics, which acts on a wave function by sending it to
the function x 7 x (x). It may be dangerous, however, to leave the domain unspecified; a
number of paradoxes arising from such sloppiness in notation are pointed out in a paper by
F. Gieres [41].
388
DEFINITION 14.30 (Extension) Let A and B be two linear operators on H.
If D A D B and if A = B for any D A , then we say that B is an

extension A or that A is the restriction of B to D A .
Remark 14.31 There is often a large choice available for the domain on which to define a
linear operator. For instance, the maximal domain on which the position operator, denoted
X , is defined is the subspace

D1 = L2 (R) ; x 7 x (x) L2 (R) .
However, this space may not be the most convenient to use. It may be simpler to restrict X to
the Schwartz space D2 = S of C class functions with all derivatives rapidly decaying. This
space is dense in H = L2 (R), and the operator X is defined on S . Moreover, this space has
the advantage of being stable with respect to X .
We must insist: the operator (X , D1 ) is not the same as the operator (X , D2 ), even if they
are both usually denoted by the same name operator X . The first is an extension of the
second (to a strictly larger space).
An important example of an operator is the momentum operator of quantum mechanics. Consider the space H = L2 [0, a] of square integrable functions on the interval [0, a]. The wave functions of a particle confined to this
interval all satisfy (0) = (a) = 0 (this is a physical constraint). We wish to
define an operator (the momentum operator P ) by
P = i h} .
If H, its derivative must be taken in the sense of distributions. But

we wish to define an operator P with image contained in H. Hence this
distribution must be a regular distribution,5 which is associated to a square
integrable function. This justifies the following definition:
DEFINITION 14.32 (Momentum operator P in L2 [0, a]) Let H = L2 [0, a].
The momentum operator P is defined, on its domain

D P = H ; H and (0) = (a) = 0 ,
by
D P
P = i h} .
Remark 14.33 (Boundary conditions for the operator P ) One may wonder what the boundary
conditions (0) = (a) really mean; as a matter of fact, , like every function of the L2 space,
is defined only up to equality almost everywhere, and therefore, it would seem that a boundary
condition is meaningless. But the fact that L2 [a, b ] implies that L1 [a, b ] and, in
particular, that is continuous.6 This explains why the boundary conditions make sense.
The momentum operator for a particle that can range over R is defined
similarly:
5
This does not imply that the function is differentiable; but it must be continuous and
its derivative exists almost everywhere.
6
Even absolutly continuous, but this is not of concern here.
Linear operators
389
DEFINITION 14.34 (Momentum operator P in L2 (R

R)) Let H = L2 (R). The
momentum operator P is defined, on its domain

D P = H ; H ,
by
D P .
P = i h}
The limit conditions () = (+) = 0 are automatically satisfied when
, H.
14.3.b Adjoint

A , D A be a linear operator on H. The
domain of A , denoted D A , is the set of vectors H for which there
exists a unique vector H satisfying
DEFINITION 14.35 (Adjoint) Let
D A
| = | A .
If D A is not empty, the operator adjoint to A, denoted A , is defined by

A = for any D A , with the preceding notation. In other words,

A = | A .
D A D A
PROPOSITION 14.36 (Existence of the adjoint) A has an adjoint if and only if
D A is dense in H.
Proof. Assume first that D A is dense. In the case = 0, the vector = 0 satisfies
the equation | = | A for all D A . So we must show that this vector is
unique. More generally, for a given arbitrary , let and satisfying the required
property. Then | = 0 for all D A . By density, we may find a sequence
(n )nN , with values in D A , which converges to . Taking the limit in the scalar
product (which is a continuous map), we obtain k k2 = | = 0
and hence = .
Conversely, assume that D A is not dense. Then the orthogonal of D A in not equal
to {0} (it is a corollary of Theorem 9.21 ). Choose any 6= 0 in this orthogonal space.
Then | = 0 for any D A . Let H. If is a candidate vector such that the
relation | = | A holds for all D A , it is obvious that the vector +
also satisfies this property. So there is no uniqueness and therefore D A = .
Remark 14.37 If D A is dense, it may still be the case that D A is equal to {0}.
In the rest of this chapter, we only consider linear operators with dense domain
in H. So the adjoint is always defined.
PROPOSITION 14.38 (Properties of the adjoint) Let A and B be linear opera-
tors.
i) A is a linear operator;
ii) (A) = A for all C;
390
iii) if A B, then B A ;
iv) if D A is dense, then A has un adjoint, denoted A , which is an extension

of A: A A .
In Section 14.4, we will determine the adjoints of some important operators.
14.3.c Bounded operators, closed operators,

closable operators
In the terminology of Hilbert spaces, a bounded operator is simply a continuous operator. This is justified by the following definition and proposition:
DEFINITION 14.39 An operator A is bounded if and only if it is bounded
on the unit ball of H, or in other words, if the set

kAk ; D A and kk 1
is bounded. If A is bounded, we denote
|||A||| = sup kAk = sup

D A
kk1
D A
6=0
kAk
.
kk
Just as in any other complete normed vector space, the following holds (it
is just a reminder of Theorem A.51 on page 583 and of sidebar 5 on page 294):
PROPOSITION 14.40 A is bounded if and only if A is continuous.
If A is bounded, it extends uniquely to a continuous linear operator defined on the

whole of H.
A property weaker than boundedness but which is also important is that

of being a closed operator.
DEFINITION 14.41 (Closed operator, graph) Let A be a linear operator. The
graph of A is the subspace

(A) = (, A) ; D A
of H H.
The operator A is closed if and only if its graph is closed,7 that is, if
for any sequence (n )nN with values in D A which converges to a limit , if
moreover the sequence (An )n converges to a vector , then we have D A
and A = .
7
For the topology defined on H H by the natural scalar product

(, )( , ) = | + | .
Linear operators
391
The reader will find the following elementary properties easy to check:
PROPOSITION 14.42 Any bounded operator is closed.
THEOREM 14.43 For any operator A, its adjoint A is closed.
DEFINITION 14.44 (Closable operator, closure) An operator A is closable if
there exists a closed extension of A; the smallest such closed extension is called
the closure of A and is denoted A.
PROPOSITION 14.45 A is closable if and only if the domain of A is dense. In
this case, we have A = (A ) and A = A .

This notion of closure of an operator may seem very abstract. We will see
below how to construct it in certain special cases important in physics.
14.3.d Discrete and continuous spectra

A , D A be an operator. An eigenvalue of A is a
complex number such that there exists a nonzero vector D A satisfying
the relation
A = .
Such a vector is an eigenvector. The set of eigenvectors of A is the discrete
spectrum of A, and is denoted d ( A).
DEFINITION 14.46 Let
Consider the momentum

operator P on the Hilbert space H = L2 (R). It is defined on D P = { H ; H} by
Example 14.47 (Discrete spectrum of the momentum operator P )
P = i h} .
Let C. We try to solve, in H, the equation
P = .
This implies
= (i/ h}) , and hence (x) = e i x/ h} , where is a constant. However,
functions of this type are square integrable only if = 0; since the zero function is not an
eigenvector, it follows that P has no eigenvalue: d (P ) = .
Example 14.48 (Discrete spectrum of the position operator X) Consider now the position op-
erator X on the Hilbert space H = L2 (R). It is defined on D X = { H ; x 7 x (x) H}

by X = , where : x 7 x (x).
Let C. As before we try to solve, in H, the equation
X = .
It follows that for almost all x R we have (x ) (x) = 0. Hence the function is zero
almost everywhere, in other words,8 = 0. Hence X has no eigenvalue: d (X ) = .
8
Since, we recall, the space L2 is constructed by identifying together functions which are
equal almost everywhere. Any function which is zero almost everywhere is therefore identified
with the zero function.
392
DEFINITION 14.49 Let A , D A
be an operator. A generalized eigenvalue

of A is a complex number such that
is not an eigenvalue of A;
there exists a sequence (n )nN of vectors of norm 1 satisfying

lim kAn n k = 0.
The set of generalized eigenvalues of A is called the pure continuous spectrum of A, and is denoted c (A). It is disjoint from the discrete spectrum:
d (A) c (A) = .
Example 14.50 (Continuous spectrum of the operator X) We saw that the operator X of the
Hilbert space L2 (R) has no eigenvalue. On the other hand, a function with graph which is a
very narrow spike around a real number should satisfy X . Hence we will check
that any real number is a generalized eigenvalue, the idea being to construct such a sequence
of spikier and spikier functions. For reasons of normalization, this is not quite the same as a
Dirac sequence.
0
Let x0 R.
p We consider the sequence (n )nN of functions defined by n (x) = n (x x0 )
and 0n (x) = n (nx); recall that the rectangle function satisfies (x) = 1 for |x| 1
and (x) = 0 otherwise.
It is easy to check that the functions n are of norm 1. For all n N, we have
Z +
Z 1/2n

1
(X x0 )n =
x 2 n (x)2 dx =
nx 2 dx =
.
12
n2
1/2n
Hence x0 is a generalized eigenvalue of the operator X .

One can also show that the continuous spectrum is a subset of R, and therefore conclude
that the continuous spectrum of X is the whole real line: c (X ) = R.
Remark 14.51 In the previous example, the sequence (n )nN converges neither in L 2 (R) nor
in the sense of distributions (because

the normalization is chosen for the purpose of the L2
p
norm). However, the sequence ( n n ) is a sequence that converges, in the sense of distributions, to T = (x x0 ), which satisfies the relation
x (x x0 ) = x0 (x x0 ),
in other words X T = x0 T in the sense of distributions. This is the (generalized) eigenvector
associated to the generalized eigenvalue!
Defining properly what is a (tempered) distribution which is an eigenvector T of an
operator A requires a proper definition of what is meant by A T . This is in fact a delicate
problem, which we will only solve for a certain type of operators (see Definition 14.68).
Example 14.52 (Continuous spectrum of the operator P ) Consider the operator P on the space
L2 (R). We have already shown that it has no eigenvalue. On the other hand, we can now show
that any p R is a generalized eigenvalue. The idea is, once more, to find a function which
looks like the function p (which suffers itself from the essential defect of not being square
integrable). For this purpose, we will simply average the functions , for values of close
to p. Such an average is, up to unimportant factors, the same as the Fourier transform of
a function similar to the rectangle function: hence it is a function that decays like 1/x at
infinity, and is therefore square integrable.
To be more precise, define
r
Z p+
h} sin(x/ h})
1
def
()
(x) d =
p (x) = p
ix
2 p
Hermitian operators; self-adjoint operators
393
for x R and > 0. It is easily checked that this function is of norm 1 in L2 (R), and that

2
(P p) ()
= Cnt 2 ,
p
so that taking = 1/n with n going to infinity shows that p is indeed a generalized eigenvalue.
Later on (see Theorem 14.82) we will see that the same operator defined on L2 [0, a] has
no generalized eigenvalue.
To finish this section, we remark that, in addition to the continuous spectrum and discrete spectrum, there exists a third type.
DEFINITION 14.53 Let A be an operator with dense domain D A . The residual spectrum of A, denoted r ( A), is the set of complex numbers such
that
/ d (A)
and
d (A ).
The spectrum of A is the union of the three types of spectra:
(A) = d (A) c (A) r (A).

The resolvant set of A is the set of complex numbers which are not in the
spectrum of A: (A) = C\(A). It is possible to prove that for any (A),
the operator A Id has a continuous inverse.
Since the residual spectrum of a self-adjoint operator is empty, we will
not mention it anymore. The interested reader can look at Exercise 14.3 on
page 403 for another definition of the various spectra, and may consult a
general reference in functional analysis, such as [71, 94].
14.4
The goal of this short section is to prove the following result, which is a
generalization of Theorem 14.7: any self-adjoint operator admits an orthogonal basis
of eigenvectors, where the basis may be a generalized basis (i.e., its elements may
be distributions). This fact is absolutely fundamental for quantum mechanics,
where an observable is defined, according to various textbooks, as being a
hermitian operator which has an orthogonal basis of eigenvectors or, better,
as a self-adjoint operator. In contrast with the case of finite dimension, the
notions of hermitian operator and of self-adjoint operator no longer coincide,
and this causes some difficulties.
Finally, we assume in this section that H = L2 (R).
Example 14.54 (A classical problem of quantum mechanics) Consider a particle in an infinite
potential sink, confined to the interval [0, a]. The operator P satisfies | P = P | for
all functions , D P . However, P has no eigenvalue, and no generalized eigenvalue. Hence
P is not an observable.
394
14.4.a
Definitions
DEFINITION 14.55 Let A be a linear operator, with dense domain D A in H.
Then A is hermitian (or symmetric) if and only if we have

, D A
A| = | A .
THEOREM 14.56 (Characterizations of a hermitian endomorphism) Let A be
an operator on H. Then the following statements are equivalent:

i) A is hermitian;
ii) D A D A and A = A for any D A ;
iii) A A .
Proof. First of all, notice that iii) is just a rephrasing of ii).

Assume that A is hermitian. Then, for any D A , the element = A satisfies
| = A| = | A ,
which proves that D A and A = A.

The converse is immediate.
DEFINITION 14.57 Let A be an operator, with D A dense in H. Then A is

self-adjoint if and only if A = A , that is, if
D A = D A
and
D A
A = A.
In a finite-dimensional space, it is equivalent for an operator to be hermitian or self-adjoint. Here, however, there is a crucial distinction: the domain
of the adjoint of a hermitian operator A may well be strictly larger than that
of A. In other words, A may be acting in the same way as A on the domain
of A, but it may also act on a larger space. Disregarding this distinction
may lead to unwelcome surprises (an apparent paradox of quantum mechanics is explained in Exercise 14.5). It is often easy to check that an operator is
hermitian, but much more delicate to see whether it is, or is not, self-adjoint.
If A is defined on the whole of H, then things simplify enormously; since
A A , it follows that the domain of A is also equal to H and so if A is
hermitian, it is in fact self-adjoint. The situation is even simpler than that:
THEOREM 14.58 (Hellinger-Toeplitz) Let A be a hermitian operator defined on H.
Then A is self-adjoint and bounded.
The adjoint of any operator is closed, and therefore we have

THEOREM 14.59 Any self-adjoint operator is closed.
This result can be made more precise:
395
PROPOSITION 14.60 Let A be a hermitian operator. Then A A ; hence D A is
dense and A is closable. Moreover, A A A , and A is self-adjoint if and only if

A = A = A .
If A is hermitian and closed, then A is self-adjoint if and only if A is hermitian.
A last result may be useful.
PROPOSITION 14.61 Let A be a hermitian operator. The defect indices of A are

the quantities, either integers or +, defined by
n+ = dim Ker(A i)
The following statements are equivalent:
and
n = dim Ker(A + i).
i) A is self-adjoint;
ii) A is closed and n+ = n = 0;
iii) Im(A i) = H.
PROPOSITION 14.62 The operator X on the space L2 [0, a] is self-adjoint and
bounded.
The operator X on the space L2 (R) is self-adjoint and unbounded.
Proof
Consider the operator X on H = L2 [0, a]. Its domain is the whole of H.
Moreover, for any , H, we have
Z +
Z +
X | =
x (x) (x) dx =
(x) x (x) dx = | X ,
X
= X and that X is self-adjoint. It is bounded because

Za
Za

x (x)2 dx a 2
(x)2 dx = a 2 kk2 ,
kX k2 =
which proves that
from which we derive in fact that |||X ||| a. (As an exercise, the reader is invited to
prove that |||X ||| = a.)
Consider now the operator X on the space L2 (R). The domain of X is

D X = H ; x 7 x (x) H .
By the same computation as before, it follows easily that X is hermitian. Lets prove
that it is even self-adjoint.
Let D X and write = X . Then, for any D X , we have X | =
| , by definition of X . Expressing the two scalar products as integrals, this gives,
for any D X ,
Z +

(x) x (x) (x) dx = 0.
In particular, taking (x) = x (x) (x) for x [M , M ] and (x) = 0 for

|x| > M , we have D X and
ZM

x (x) (x)2 dx = 0,
M
396
hence (x) = x (x) almost everywhere on [M , M ], and hence almost everywhere

on R since M was arbitrarily large. It follows that x 7 x (x) L2 (R), which means
D X , and in addition = X .
Thus, we have shown that D X D X and that X is the restriction of X , that is,
X X . Since we know that X X , we have X = X .
Remark 14.63 The operator X on L 2 [0, +[ is also self-adjoint, the previous proof being
valid without change.
PROPOSITION 14.64 The operator P on L2 [0, a] is hermitian (symmetric), but
not self-adjoint.
Proof. A function L2 [0, a] belongs to the domain of P if and only if it satisfies

the following two conditions:
and are in L2 [0, a],
(0) = (a) = 0.
()
()
(See Remark 14.33 on page 388 for the meaning of these boundary conditions.) The
subspace D P is dense. Moreover, for any , D P , we have
Za
i h} (x) (x) dx
P | =
0
Za

= i h}
(x) (x) dx + i h} (a) (a) (0) (0) = | P ,
0
since the boundary terms cancel out from condition (). This establishes that P is
symmetric.
However the domain of P is strictly larger than that of P . Indeed, any function
which satisfies only the condition () is such that P | = | P for any D P .
Using techniques similar to that in the proof of Theorem 14.62, it is possible to prove
the converse inclusion, namely, that

D P = ; , L2 [0, a] .
On D P , the operator P is defined by P = i h}. So the adjoint of P acts in the
same manner as P on the domain of P , but it has a larger domain.
14.4.b Eigenvectors
The first important result concerning hermitian operators is in perfect analogy with one in finite dimensions: eigenvectors corresponding to distinct
eigenvalues are orthogonal.
THEOREM 14.65 Let A be a hermitian operator. Then:
i) For any D A , | A is real.
ii) The eigenvalues of A are real.
iii) Eigenvectors of A associated to distinct eigenvalues are orthogonal.

iv) If (n )nN is an orthonormal Hilbert basis where each n is an eigenvector
of A, then the discrete spectrum of A is exactly {n ; n N}; in other words,
there is no other eigenvalue of A.
397
Remark 14.66 Concerning point iv), there is nothing a priori that ensures that such a basis of
eigenvectors exists.
PROPOSITION 14.67 The residual spectrum of a self-adjoint operator A is empty.

Proof. Let r (A). Then d (A ) = d (A) since A = A ; but the eigenvalues
of A are real, hence = and therefore d (A), which contradicts the definition
of the residual spectrum.
14.4.c Generalized eigenvectors

If A is an operator and is a generalized eigenvalue, we are going to try to
solve the equation A T = T in the sense of (tempered) distributions. The
first thing to do is to define A T . We have already seen the method: start by
looking at what happens in the case of a regular distribution T associated
with a function S , that is, defined by
Z +
T , = | =
(x) (x) dx.
If A is a symmetric operator, we have

S
hence

A = | A ,

A T = T A .
This property can be extended, as long as A is defined for all S and is

also in S :
DEFINITION 14.68 (Image of a distribution by a symmetric operator) Let A
be a symmetric operator, such that

S DA ;
S is stable under A;
A is continuous on S (for its topology).
Then A is of class S . For any tempered distribution T S , the tempered

distribution A T is defined by
S
A T | = T | A .
Exercise 14.1 Prove that A T is indeed a tempered distribution.
DEFINITION 14.69 (Generalized eigenvector) Let A be a symmetric operator

of class S . A nonzero tempered distribution T is an eigendistribution or a
generalized eigenvector, associated to the eigenvalue , if and only if
A T = T ,
in other words if T | A = T | for any S .
398
The main result is the following:

THEOREM 14.70 (Spectral theorem) For any self-adjoint operator of a Hilbert
space H which is of class S , there exists an orthonormal basis of eigendistributions

which are associated to real eigenvalues or generalized eigenvalues.
Remark 14.71 There exists a much more general version of this theorem, which applies to any
self-adjoint operator. This generalized spectral theorem states that there always exists a spectral
family, but this notion is beyond the scope of this book. It is discussed in many books, for
instance [71, III.5, 6], [95, p. 424], and [38, vol. 4, ch.1, 4].
2
PROPOSITION
14.72 The

momentum operator P on the space L (R) is of class S .
The family p| ; p R is a generalized orthonormal basis of generalized eigenvectors for P .
Proof. We already know, by Theorem 14.25, that this family is a generalized basis.
The orthonormality is the statement that
p, p R
p| p = (p p ).
We now check that, denoting by p the regular distribution associated to the function
p (x) = (2 h})1/2 e i x/ h} , we have P p = p p .
Indeed, let S . Then we have

P p = p P
by Definition 14.68
f

= p i h} = i h} (p)
= p (p)
e
since F [ ] (p) = (i p/ h}) (p)
e

= p p = p p
since p is real.
()
Remark 14.73 The distribution p is, up to normalization, the limit of the p defined in
Example 14.52.
2
PROPOSITION
14.74

The operator X on L (R) is of class S . Moreover, the
family x| ; x R is a generalized orthonormal basis of generalized eigenvectors

for X .
Proof. It is only necessary to check that X (x x0 ) = x0 (x x0 ) in the sense of

Definition 14.68. For any S , we have indeed

X (x x0 ) = (x x0 ) X = (X )(x0 ) = x0 (x0 ) = x0 (x x0 ) .
14.4.d Matrix representation
Let A be an operator on H, with dense domain as usual. Let (n )nN be an

orthonormal basis of H, all elements of which are also in D A . The associated
matrix elements are defined by
def
a i,k = i | A k .
The question is: is it true or not that the knowledge of the infinite matrix
M = (a ik ) i,k is sufficient to reconstruct entirely the operator A?
399
DEFINITION 14.75 The operator A has a matrix representation associated

with ( n ) n if and only if

P
P
P
=
xk k
A =
a ik xk i .
(14.1)
i=0
k=0
k=0
Assume that A has a matrix representation. Using Diracs notation, we

have xk = k | ; still following Dirac, define9
def
i | A |k = i | Ak
and derive
A | =
X
i=0 k=0
| i i | A |k k | ,
which is again the closure relation. Symbolically, this is written:

X
A=
a ik | i k | .
i,k
The difficulty is that these formulas are not always correct, since the righthand side of (14.1) is not always correctly defined for all D A . It is easy
to see that, if A is bounded, then by continuity everything works out. More
generally:
THEOREM 14.76 (Matrix representation) If A be a self-adjoint operator, then it
has a matrix representation associated to any orthonormal basis of D A .

This property is also true if A is bounded, or if it is hermitian and closed.10
As an example, the operators of orthogonal projection on a closed subspace

are self-adjoint operators which admit a matrix representation:
THEOREM 14.77 Let (k )kN be an orthonormal family, and let M be the closure
of Vect{k ; k N}. Then M is a closed subspace. Let M be its orthogonal

complement. Any vector H can be decomposed in a unique way as
= M +
with M M and M .
This is where the trap lies for those who manipulate the notation blindly. The definition
def
| A = A |, which is often found, permits two readings of the expression | A |:

as A which imposes that | D A and | H;

as | A which imposes that | H and | D A .
Those are not quite equivalent always because of the issue of the exact domain of an operator.
Note that, for a self-adjoint operator, there is no ambiguity; for a hermitian operator, mistakes
can arise because of this [41].
P
P
10
2
For instance, if
i=0 |a ik | < + for all k N, the operator A : k 7
i=0 a ik i
defines, by density of Vect{ i ; i N}, a closed operator, which is hermitian if aki = a ik .
400
The operator which associates M to is the orthogonal projector on M. It

satisfies 2 = and = . In matrix representation it is given by
=
k=0
|k k | .
Let A be a self-adjoint operator, and assume that there is an orthonormal

Hilbert basis (n )nN of eigenvectors of A, with corresponding eigenvalues
(n )nN (recall that this is by no means always the case). Then
A=
n=0
that is,
D A
n |n n | ,
A | =
and
X
n=0
n n | |n ,
X
2
2n n | converges.
D A if and only if
With this the formal calculus of the operator A can be defined:
DEFINITION 14.78 (Function of an operator) Let A be a self-adjoint operator,
(n )nN an orthonormal Hilbert basis of eigenvectors of A, associated to the

eigenvalues (n )nN . Let F : R R be a measurable function. The operator
F (A) is defined by
P
F (A) =
F () |n n | ,
n=0
with domain
n
o
X
D F (A) = H ;
F ()2 |n | |2 < + .
n=0
All the n are in D F (A) , hence D F (A) is also dense. In fact we have [95, 5.8]:
THEOREM 14.79 The operator F (A) is self-adjoint.
Since the identity operator is obviously self-adjoint and defined on H,

for any orthonormal Hilbert basis (n )nN and any sequence (n )nN of real
numbers, the operator
X
def
B=
n |n n |
n=0
is a self-adjoint operator.
This technique can provide the extension of an operator which is merely
hermitian, but admits an orthonormal basis of eigenvectors, to a self-adjoint
operator defined therefore on a larger domain, and which is the only one
for which many computations make sense.
401
For instance, consider the operator H0 defined on D H0 = D(R3 ) (space of

C functions with bounded suppoert), defined by
h}2
D
H0 (x) =
+V (x) (x),
2m
where V is a given function. In many cases, H0 is essentially self-adjoint,
that is, the closure H0 is self-adjoint; this closure of H0 is then given by
H0 = H0 . The self-adjoint extension H = H0 , is the hamiltonian operator
of the system. It is this operator which is observable. In practice, it is
obtained by finding an orthonormal basis (n )nN for H0 and applying the
method above.
Example 14.80 (Hamiltonian for the harmonic oscillator) Consider a harmonic oscillator with
mass m and pulsation . We first define a formal hamiltonian H0 on D H0 = S by

P2
m2 2
| +
X | .
2m
2
This operator is hermitian, and has a Hilbert basis (n )nN such that
1
h}.
H0 |n = En |n
with En = n +
2
The functions n are of the type exp(m2 x 2 /2 h}) hn (x), where hn is a Hermite polynomial
of degree n [20, Ch. V C.2]. The operator H0 is not self-adjoint.
Let us now define the operator H by
P
P
H=
En |n n | ,
i.e.,
H | =
En n | |n
H0 | =
n=0
n=0

2
for any function | such as the series En2 n | converges.
The operator H is self-adjoint; it is an extension of H0 , which is called the hamiltonian
of the harmonic oscillator.
P
An example of this is given in Exercise 14.5.
14.4.e Summary of properties of the operators P and X

R)) The operator P on L2 (R) is selfPROPOSITION 14.81 (Operator P on L2 (R
adjoint (and hence closed). It is defined on
D P = { H ; H}.
The discrete spectrum

of P is empty; its continuous spectrum is equal to R. The family

p| ; p R is a generalized basis of eigenvectors such that

p p = (p p ),
(orthonormal basis)
Z +
Id =
| p p| d p.
(closure relation)
Proof. See Example 14.52 for the continuous spectrum and Theorem 14.72 for the
basis of eigenvectors.
402
PROPOSITION 14.82 (Operator P on L2 [0, a])
The operator P on L2 [0, a] is hermitian (symmetric), but not self-adjoint. Its domain
is

D P = H ; H and (0) = (a) = 0 .
P is closed (and therefore admits matrix representations).

The discrete spectrum of P is empty; its continuous spectrum is also empty and its
residual spectrum is equal to C.
Proof. The first property has already been proved in Proposition 14.64.
Before trying to prove the absence of any eigenvalue, even generalized, we can try
to reason physically. A vector such that P is close to is characterized by a
very small uncertainty p with respect to the momentum, which, according to
Heisenbergs uncertainty principle is only possible if x is very large. But, here, x
is bounded by a.
Such an argument is not simply heuristic: it can be justified by noting that P
and X (on L2 (R)) are unitarily equivalent via the Fourier transform, namely, we have
P = F X F 1 . Let now p0 be any real number; we will show that p0 is not a generalized
eigenvalue. Let be a function with norm 1. The Fourier transform e of (P p0 ) =
i h} p0 is p 7 (p p0 ) (p).
e
The Parseval-Plancherel identify shows that

2
(P p0 ) 2 = (p p0 ) (p)
e .
2
The norm squared e is simply the average value, in the state corresponding to ,

of the quantity (p p0 )2 , which was denoted (p p0 )2 in the section on uncertainty

relations. It can be checked that the quantity (p p0 )2 is smallest for p0 = p,
hence

e 2 =
(p p )2 (p)2 .
(p)
0
By the uncertainty relations (Theorem 13.17), we know that p h}/2x h}/2a,
which finally implies
2
2
e h} .
4a
2

2
Consequently, it is not possible to make p 7 (p p0 ) (p)
e = (P p0 ) as
small as possible, so that p0 is not a generalized eigenvalue.
One easily checks that the domain of P is D P = { H ; H}, and that any
complex is an eigenvalue of P , with associated eigenvector : x 7 e i x/ h} . This
proves that r (P ) = C.
R)) The operator X on L2 (R) is selfPROPOSITION 14.83 (Operator X on L2 (R
adjoint; its domain is
D X = { L2 (R) ; x 7 x (x) L2 (R)}.
The discrete spectrum of X is empty; its continuous spectrum is R. The family

x| ; x R
is a generalized basis of generalized eigenvectors such that

x x = (x x ),
Z +
Id =
|x x| dx.
(orthonormal basis)
(closure relation)
Exercises
403
Proof. See Theorem 14.62 to show that X is self-adjoint. See Remark 14.50 for the
spectrum.
PROPOSITION 14.84 (Operator X on L2 [0, a]) The operator X on L2 [0, a] is

self-adjoint and bounded. Its domain is D X = L2 [0, a]. The discrete spectrum of X
is empty. Its continuous spectrum is [0, a]. The family

x| ; x [0, a]
is a generalized basis of eigenvectors of X on L2 [0, a].
Proof. See Theorem 14.62 to show that X is self-adjoint. The computation of the
spectrum is left as an exercise; the proof in Remark 14.50 can be used as a clue.
EXERCISES
Exercise 14.2 Check that any eigenvector is also an eigendistribution.
Exercise 14.3 Let A be a linerar operator defined on the subspace D A = W . For any
complex number , the operator A = A Id is also defined on W . One and only one of
the following four cases applies for (see, e.g., [94]):
A is not injective; then is an eigenvalue, or equivalently belongs to the discrete

spectrum.
A is injective, Im(A ) is dense but the inverse11 of A is not continuous; then is a
generalized eigenvalue and belongs to the continous spectrum.
A is injective but Im(A ) is not dense; then belongs to the residual spectrum;
A is injective, Im(A ) is dense and A1 is continuous (and hence may be extended by

continuity to an operator defined on the whole of H); then belongs to the resolvant
set of A.
Show that those definitions agree with those in Definitions 14.46 and 14.49.
Note that in some books (for instance [4]), the definitions somewhat different, and the
discrete and continuous spectrum, for instance, are not necessarily disjoint. (This may actually
be desirable to account for natural situations, such as some hyperbolic billiards arising in
arithmetic quantum chaos, where there is a seemingly chaotic discrete spectrum embedded
in a very regular continuous spectrum; see, e.g., [78].)
Exercise 14.4 (Continuous spectrum of X) Consider the Hilbert space L 2 [a, b ]. Show that
X has no eigenvalue, but that any [a, b ] is a generalized eigenvalue, that is, d (X ) =
and c (X ) = [a, b ].
Exercise 14.5 (A paradox of quantum mechanics, following F. Gieres [41])
Consider a particle, in an infinite potential sink, constrained to remain in the interval [a, a].
Physical considerations impose that the wave functions satisfy (a) = (a) = 0. The
hamiltonian is the operator on L2 [a, a] defined by H = h}2 /2m. Assume H is a
self-adjoint operator.
11
The target space of A can be restricted to Im(A ); then A is a bijection from W

to Im(A ).
404
i) Solve the eigenvalue equation H = E , associated to the boundary conditions

(a) = 0; deduce from this the spectrum (En )nN and an orthonormal basis of
eigenfunctions (n )nN .
ii) Consider the wave function given by (x) = (a 2 x 2 ), where is a normalizing
factor so that | =
1. Show that H | H > 0. Isnt it the case however that
H 2 = 0 and hence H 2 = 0? How should the average value of the energy
square in the state be computed? Comment.
Exercise 14.6 (Uncertainty principle) Let A and B be two self-adjoint operators on H (i.e.,
observables). Let H be a normed vector (i.e., a physical state) such that
DA DB ,
A D B ,
B D A .
and
The average value of A and the uncertainty on A in the state are defined by
def
a = | A
Show that
and
a = kA ak .

a b 21 [A , B] ,
def
where [A , B] = A B B A.
Check that [X , P ] = i h} Id and deduce from this that x p h}/2.
SOLUTIONS
i) Solving the eigenvalue equation leads to the discrete spectrum given by:
En =
h}2 2 n2
8ma 2
(n 1)
associated with the eigenfunctions

p
n x
2mEn x
1
1
n (x) = p cos
= p cos
a
h}
a
2a
p

2mEn x
1
1
n x
n (x) = p sin
= p sin
a
h}
a
2a
if n is odd,
if n is even.
It is easy to check that the sequence (n )nN is a complete orthonormal system of

eigenfunctions of H .
p
ii) First compute that = 15/4a 5/2 . Then another computation shows that
kH k2 = H | H =
15 h}4
.
8m 2 a 4
2
Moreover, since
fourth
the
derivative of is zero, it is tempting to write H = 0 and,
2

consequently, H = 0; which is quite annoying since H is self-adjoint, and one
would expect to have identical results.
The point is that one must be more careful with the definition of the operator H 2 .
Indeed, the domain of H is
D H = { H ; H and (a) = 0};
405
hence H
/ D H and therefore it is not permitted to compute H 2 as the composition
H (H ). (It may be checked that the elements of D H 2 must satisfy (a) = 0.) On
the other hand, using the orthonormal Hilbert basis (n )nN , we can write
H=
En |n n |
En2 |n n | .
n=1
and, using the techniques of Section 14.4.d, Definition 14.78, it is possible to define a
new operator, denoted H 2 , by
def
H2 =
n=1
Notice that since the eigenfunctions n are in the domain of H 2 , we have

n 1
H 2 n = H 2 n .
This new operator is thus an extension

of the previous one. Its domain is the subspace
P 4
2
of functions such that
En n | < +. Then one checks that the function
, which did not belong to the domain12 of H 2 , is indeed in the domain of H 2 .
This extension is self-adjoint. A (boring) computation with Fourier series shows that
n 1
2 def X
15 h}4
En2 | n n | =
H =
,
8m 2 a 4
n=1
2
i.e., that H = H | H as it should be.
Solution of exercise 14.6. First of all, notice that a and b are real numbers since A and B
are hermitian operators (Theorem 14.65). Also, notice that

(A B B A) = | A B | B A = A| B B| A = 2Im A| B .
Since a and b are real, we can expand

(A a)(B b ) = A| B ab b a + ab kk2 = A| B ab ,
and obtain

(A B B A) = 2 Im (A a)(B b ) .
The conclusion now follows using the Cauchy-Schwarz inequality:

a b (A a)(B b ) = Im (A a)(B b )

1

Im A| B = [A , B] .
2
12
This is of course linked to the Fourier series of , which converges uniformly but cannot
be differentiated twice termwise.
Chapter
15
Green functions
This chapter, properly speaking, is not a course on Green functions, and it does not
really introduce any new object or concept. Rather, through some simple physical
examples, it explains how the various techniques discussed previously (Fourier and
Laplace transforms, conformal maps, convolution, differentiation in the sense of
distributions) can be used to solve easily certain physical problems related to linear
differential equations.
The first problem concerns the propagation of electromagnetic waves in the vacuum. There the Green function of the dAlembertian is recovered from scratch, as
well as the retarded potentials formula for an arbitrary source distribution.
The second problem is the resolution of the heat equation, using either the Fourier
transform, or the Laplace transform.
Finally, it will be seen how Green functions occur naturally in quantum mechanics.
15.1
Generalities about Green functions
Consider a linear system with input signal I and output (or response)
signal R. This is described by an equation of the type
(R) = I .
What we called input signal and output signal could be many things, such
as:
electrical signals in a circuit (for instance, power as input and response
of a component as output);
408
Green functions
charges and currents as input and electromagnetic fields as output;
heat sources as input and temperature as output;
forces as input and position or velocity as output.
The operator is linear and continuous. It may depend on variables such

as time or position. In this chapter, we are interested in the case where it is
translation (spacial or temporal) invariant.1 Most of the time, is a differential
operator (such as a Laplacian, dAlembertian, etc.).
Since is at the same time continuous, linear, and translation invariant,
it is known that it may be expressed as a convolution operator; hence there
exists2 a distribution D such that I = (R) = D R. The Green function of
the system is defined to be any distribution G satisfying the equation
(G ) = D G = G D = ,
where the Dirac distribution applies to the spacial or temporal variable

or variables of the system. Computing the convolution product of the input
signal with the Green function, we get
G I = G (D R) = (G D ) R = R = R,
which shows that knowing the Green function is sufficient in principle to

compute R from the knowledge of the input I . However, complications arise
when different Green functions exist in different algebras (for instance, in
electromagnetism, two distinct Green functions lead, respectively, to retarded
and advanced potentials).
Thus, this chapters goal is to provide some examples of explicit computation techniques available to compute the Green functions, using mainly the
Fourier and Laplace transforms. The strategy used is always the same:
i) take the Fourier or Laplace transform of the equation;
ii) solve the resulting equation, which is now algebraic;
iii) take the inverse transform.
This may also be symbolized as follows:
e]
e Ge = 1
e
G = F [1/ D
D
Ge = 1/ D
or
or
or
D G = =
=
=
b
b Gb = 1
b
G 1/ D.
D
Gb = 1/ D
Despite the simplicity of this outline, difficulties sometimes obstruct the way;
we will see how to solve them.
1
2
Still, an example which is not translation invariant is treated in Section 6.2.a on page 165.
It suffices to put D = ().
A pedagogical example: the harmonic oscillator
409
15.2
We treat here in detail a simple example, which contains the essential difficulties concerning the computation of Green functions, without extraneous
complications (temporal and spacial variables together).
A Green function of the harmonic oscillator is a function t 7 G (t) that
is a solution, in the sense of distributions, of the equation
G (t) + 02 G (t) = (t),
(15.1)
where 0 > 0 is a fixed real number. This amounts to saying that G is a

convolution inverse of the operator ( + 02 ), and we have already seen
expressions for this in Theorem 8.33 on page 239.
We discuss here two methods that lead to the discovery of such a Green
function of the harmonic oscillator (in addition to the manual method of
variation of the constant,3 which does not generalize very well).
3
We quickly recall how the differential equation

y + y = f (t)
(E )
can be solved, where f (the exterior excitation) is a continuous function, and where we put
0 = 1. This is done in two steps:
Resolution of the homogeneous equation. We know that the space of solutions of the
associated homogeneous equation (namely y + y = 0) is a two-dimensional vector space, a
basis of which is, for instance, (u, v) = (cos, sin).
Search for a particular solution. We use the method of variation of the constants,
generalized for order 2 equations, as follows: look for a solution of the type y0 = u + v,
where and are functions to be determined. Denoting by Wu,v the Wronskian matrix
associated to the basis (u, v), which is defined by
u v
Wu,v =
,
u v
those functions are given by

0
cos t
1
=
W
=
u,v
f
sin t
that is,
sin t
cos t
,
f (t)
= sin(t) f (t)
and
= cos(t) f (t),
so that (up to a constant, which can be absorbed in the solution of the homogeneous equation)
we have
Zt
Zt
and
(t) =
cos(s) f (s) ds.
(t) = sin(s) f (s) ds
0
Using the trigonometric formula sin(t s) = sin t cos s cos t sin s, we deduce
Zt
y0 (t) =
f (s) sin(t s) ds.
0
410
Green functions
15.2.a Using the Laplace transform

Denote by G (p) the Laplace transform of G (t). In the Laplace world, equation (15.1) becomes
(p 2 + 02 ) G (p) = 1.
Since G is a function, and not a distribution, we therefore have
1
G (p) = 2
p + 02
and hence, using for instance a table of Laplace transforms (see page 613), we
have
1
H (t) sin(0 t).
G (t) =
0
Remark 15.1 On the one hand, this is quite nice, since the computation was very quick. On
the other hand, it may be a little disappointing: we found a single expression for a Green
function. More precisely, we found the unique inverse image of + in D+ , and not the
one in D (see Theorem 8.33 on page 239). This is quite normal: by definition of the unilateral
Laplace transform, only causal input functions are considered (hence they are in D + ).
15.2.b Using the Fourier transform

Denote by Ge() the Fourier transform of G (t). Equation (15.1) becomes
(42 2 + 02 ) Ge() = 1
(15.2)
in the Fourier world. A first trial for a solution, rather naive, leads to
1
Ge() = 2
.
(naive solution)
0 42 2
Then we should take the inverse Fourier transform of Ge().

But here a serious difficulty arises: the function Ge as defined above has
two real poles at = 0 /2, and hence is not integrable (nor is it square
integrable). So the inverse Fourier transform is not defined.
In fact, the reasoning was much too fast. The solutions of equation (15.2)
in the sense of distributions are those distributions of the form

1
0
0
Ge() = pv 2
+
+
,
(15.3)
2
2
0 42 2
Thus, the general solution of the equation are given by
Zt
y(t) =
f (s) sin(t s) ds + cos t + sin t .
|
{z
}
0
free part
We recognize, for t > 0, the convolution of the function G : t 7 H (t) sin t and the truncated
excitation function H (t) f (t). The function G is the causal Green function associated to the
equation of the harmonic oscillator (see Theorem 8.33 on page 239).
411
according to Theorem 8.11 on page 230, where we denote

1
1
1
1
pv 2
=
pv
.
pv
0 42 2
2
2 + 0
2 0
The last two terms of equation (15.3) are the Fourier transforms of the functions t 7 e i0 t and t 7 e i0 t , in which we recognize free solutions
(those that satisfy y + y = 0). We could leave them aside, but it is more
useful to find interesting valus of and .
Indeed, if we take = = 0, this amounts to use as integration contour
in the inverse Fourier transform the following truncated real line (where 0 =
0 /2):
2
0
in the limit 0.
As seen in Section 8.1.c on page 225, we may add to 0 a positive or
negative imaginary part i, which amounts to closing the contour by small
half-circles either above or under the poles. Replacing each occurence of 0 by
0 + i, or in other words, putting

1
1
1
e
G () =
,
4 0 + 0 + i 0 i
leads to a complex-valued function G , which is a priori uninteresting. We will

rather define

1
1
1
def
(ret)
e
,
G () = lim+
0 4 0 + 0 i
0 i
or, equivalently, put = = 1/40 :

1
1
1
pv
+ i ( + 0 ) pv
i ( 0 ) .
Ge(ret) () =
4 0
+ 0
0
In the inverse Fourier transform formula
Z
G (ret) (t) = Ge(ret) () e 2it d,
this corresponds to the integration contour sketched below:

given by the path
G (ret) (t)
Green functions
412

1
Since the inverse Fourier transform of 2i
pv 1 + 2 is H (t), it follows that

1
1
1
T.F.1
+ ( + 0 ) H (t) e i0 t
pv
2i
+ 0
2
and

1
1
1
T.F.1
+ ( 0 ) H (t) e i0 t .
pv
2i
0
2
Hence the inverse Fourier transform of Ge(ret) is

1
1
H (t) i e i0 t ie i0 t =
H (t) sin(0 t).
G (ret) (t) =
20
0
Remark 15.2 It is no more difficult, if the formula is forgotten and no table of Fourier
transforms is handy, to perform the integration along the contour indicated using the residue
theorem. The sign of t must then be watched carefully, since it dictates whether one should
close the integration contour from above or from below in the complex plane.
In the case discussed here, for t < 0, the imaginary part of must be negative in order
that the exponential function e 2i t decay fast enough to apply Jordans second lemma. In
this case, the integrand has no pole inside the contour, and the residue theorem leads to the
conclusion that G (ret) (t) = 0. If, on the other hand, we have t > 0, then we integrate from
above, and the residue at each pole must be computed; this is left as an exercise for the reader,
who will be able to check that this leads to the same result as previously stated.
In the following, the integration will be executed systematically using the method of
residues.
Similarly, we will define

def
Ge(ad) () = lim+
1
4 0
or, equivalently
Ge(ad) () =

1
1
,
+ 0 + i 0 + i

1
1
1
pv
i ( + 0 ) pv
+ i ( 0 ) ,
4 0
+ 0
0
which now corresponds to the integration contour
given by the path

G (ad) (t)
As before we compute the inverse Fourier transform of Ge(ad) , and get

G (ret) (t) =
1
H (t) sin(0 t)
0
and
G (ad) (t) =
1
H (t) sin(0 t).
0
413
If we chose to take = = 0, we obtain the average of the two preceding

Green functions, namely,

1
Ge(s) = Ge(ret) + Ge(ad)
2
(where s stands for symmetric). Looking separately at t > 0 and t < 0,

this has a more compact expression:
G (s) (t) =

1
sin 0 |t| .
2
So we have three distinct possibilities for a Green function of the harmonic

oscillator the reader is invited to check by hand that all three do satisfy
the equation G + 02 G = . The function G (ret) has support bounded
on the left it is the causal, or retarded, Green function. G (ad) has support
bounded on the right, and is the advanced Green function. As for G (s) , its
support is unbounded and belongs to no convolution algebra (in particular,
its convolution product with itself is not defined), which explains why it did
not come out in earlier chapters, although it is no less interesting than the
others.
The retarded Green function is the most useful. Indeed, if a harmonic
oscillator which is initally at rest (for t = t0 ), with pulsation 0 , is submitted
to an excitation f (t) which is zero for t < t0 (the function f is then an
element of the convolution algebra D+ ), the length y of the oscillator satisfies
y (t) + y(t) = f (t),
t R
the solution of which, taking initial conditions into account, is

y = G (ret) f ,
that is,
y(t) =
t
t0
f (s) G (ret) (t s) ds.
Remark 15.3 Other complex-valued Green functions can also be defined. Those are of limited
interest as far as classical mechanics is concerned, by are very useful in quantum field theory
(see the Feynman propagator on page 431).
Green functions
414
15.3
Electromagnetism and the dAlembertian operator
DEFINITION 15.4 The dAlembertian operator or simply dAlembertian is
the differential operator
1 2
.
c 2 t2
(Certain authors use a different sign convention.)
def
15.3.a
Computation of the advanced and retarded

Green functions
We assume known the charge distribution ( r , t) and the current distribution

j( r , t), and we want to find the electromagnetic fields produced by those
distributions, up to a free solution. The easiest method is to start by looking
for the potentials ( r, t) and A( r , t) which, in Lorenz gauge, are related to the
sources by the following second order partial differential equation:

/0
=
.
(15.4)
A
0 j
Hence we start by finding the Green function associated to the dAlembertian,
that is, a the distribution G ( r , t) satisfying the same differential equation
with a source ( x) (t) which is a Dirac distribution, in both space and time
aspects:
G ( r , t) = ( r ) (t)
(15.5)
(Green function of the dAlembertian).

def
def
Put A( r , t) = ( r, t), A( r, t) and j( r, t) = ( r , t), j ( r , t) . The general
solution of (15.4) is given by
ZZZ Z +
A( r , t) = [G j]( r , t) =
G ( r r , t t ) j( r , t ) dt d3 r
R3
(up to a free solution).

Passing to the Fourier transform in equation (15.3.a) with the convention4
ZZZZ
def
G ( k, z) =
G ( r , t) e i( k rz t) d3 r dt
for all k R3 and z R,
4
Notice the sign in the exponential; this convention is customary in physics.
415
An illegitimate son abandoned in front of the Saint-Jean-le-Rond

Church, the young Jean le Rond dAlembert (17171783) studied
in Jesuit schools and directed his attention to mathematics. Fond
of success in the salons, he was interested in Newtons theories,
then much in vogue. He collaborated with Diderot to edit and
write the Encyclopdie; in particular, writing the Discours prliminaire and many scientific articles. He studied complex analysis,
the equation of vibrating strings, and probability theory, and gave
the first essentially complete proof of the fundamental theorem
of algebra (any polynomial with real coeffecients can be factorized into polynomials of degree 1 or 2), which was completely
proved by Gauss.
we obtain the relation5

z2
2
k
G ( k, z) = 1.
c2
A possible solution to this equation is6

G ( k, z) =
c2
,
z2
c 2 k2
which leads, after performing the inverse Fourier transform, to

ZZZZ
1
c2
G ( r , t) =
e i( k rz t) d3 k dz.
(2)4
c 2 k2 z2
(15.6)
However, this integral is not well defined, because the integrand has a pole.
To see what happens, we start by simplifying the expression by means of an
integration in spherical coordinates where the angular components can be
easily dealt with. So we use spherical coordinates on the k-space, putting
d3 k = k 2 sin dk d d, with the polar axis in the direction of r . We have
5
To avoid sign mistakes or missing 2 factors, etc, use the following method: start from
the relation G = , multiply by e i ( k r z t) , then integrate. Then a further integration by
parts leads to the equation above.
6
We know that there exist distributions (free solutions) satisfying
2

z
2
k
T ( k, z) = 0,
c2
as, for instance, T ( k, z) = 3 ( k) (z) or T ( k, z) = ( k, z) (z 2 k 2 c 2 ). Adding these to the

general solution changes in particular the behavior (decay) at infinity; in general, the solution
will decay at infinity if its Fourier transform is sufficiently regular.
Green functions
416
therefore
ZZZZ
1
c2
e i(kr cos z t) k 2 sin dk d d dz
2
2
2
(2)4
c k z
Z +
Z +
c2
k
i
z
t
=
sin(kr)
e
dz dk.
2 2
2
43 r 0
c k z
G ( r , t) =
For the evaluation of the integral in parentheses, there is the problem of the
pole at z = ck. To solve this difficulty, we move to the complex plane
(that is, we consider that z C), and avoid the poles by deforming the
contour of integration. Different choices are possible, leading to different
Green functions.
Remark 15.5 It is important to realize that, with the method followed up to now,
there is an ambiguity in (15.6). The ways of deforming the contour lead to radically
different Green functions,7 since they are defined by different integrals. These functions
differ by a free solution (such that f = 0). The solution is the same as in the
preceding section concerning the harmonic oscillator: take the principal value of the fraction at
issue, or this principal value with an added Dirac distribution i at each pole.
We will make two choices, leading to functions which (for reasons that
will be clear at the end of the computation) will be named G (ad) ( r , t) and
G (ret) ( r , t):
given by the path
G (ad) ( r , t)
ck
ck
and
given by the path

G (ret) ( r , t)
ck
ck
First, we compute the retarded Green function. For negative values of the
time t, the complex exponential e i z t only decays at infinity if the imaginary
part of z is positive; hence, to be able to apply the second Jordan lemma, we
must close the contour in the upper half-plane (see page 120):
R +
ck
7
ck
Radically different for a physicist. Some may be causal while some others are not, for
instance.
417
Since the poles are located outside this contour, the integral is equal to
zero:
G (ret) ( r , t) 0
for t < 0.
For positive values of t, we have to close the contour in the lower half-plane,
and both poles are then inside the contour. Since the path is taken with the
clockwise (negative) orientation, we find that for t > 0 we have
i z t

Z
e
e i z t
dz
=
2i
Res
;
z
=
ck
2 2
2
c 2k 2 z 2
(ret) c k z

i i ckt
=
e
e i ckt ,
ck
and hence
Z +

c2
ik i ckt
i ckt dk
sin(kr)
e
e
43 r 0
ck
Z +
c
=
e ik(r c t) e ik(r c t) e ik(r +c t) + e ik(r +c t) dk
2
8 r 0
Z +
c
=
e ik(r c t) e ik(r +c t) dk
2
8 r

c
=
(r c t) (r + c t) ,
4r
R ik x
since e dk = 2 (x) (see the table page 612). Because t is now strictly
positive and r 0, the distribution (r + c t) is identically zero on the whole
space. In conclusion, we have obtained
G (ret) ( r , t) =
G (ret) ( r , t) =
c
(r c t)
4r
for t > 0.
Note that this formula is in fact also valid for t < 0, since the remaining
Dirac distribution is then identically zero.
Similarly, the advanced Green function is zero for positive times and is
given for all t by
c
(r + c t) for t < 0,
4r
G (ad) ( r , t) 0
for t > 0.
G (ad) ( r , t) =
Physical interpretation
The retarded Green function is represented by a spherical shell, emitted at
t = 0 and with increasing radius r = c t. The fact that this pulse is so localized
explains why a light flash (from a camera, for instance) is only seen for a time
Green functions
418
equal to the length of time it was emitted8 but with a delay proportional to
the distance. Moreover, the amplitude of the Green function varies like 1/r,
which is expected for a potential.
15.3.b Retarded potentials

Coming back to equation (15.4) for potentials in Lorenz gauge, we obtain using
the retarded Green function for the dAlembertian that

(ret)
(ret)
( r , t),
A (ret) ( r , t) = G (ret) 0 j ( r , t),
( r , t) = G
0
or
(ret)
c
( r , t) =
40
ZZZ Z
R3

( r , t )
k r r k c (t t ) dt
kr r k
d3 r .
We can integrate9 the time variable t . Since c (r c t) = (t r/c ), only

the value t = t k r r k/c need be kept, which gives

ZZZ
r , t k r r k/c 3
1
(ret)
( r, t) =
dr.
40
k r r k
R3
Each integration element therefore gives the contribution of the charge element situated at the point r at the retarded time = t k r r k/c .
Similarly, we find that

ZZZ
, t k r r k/c
j
r
0
A(ret) ( r , t) =
d3 r .
4
k r r k
R3
Remark 15.6 Only a special solution has been described here. Obviously, the potential in
space is the sum of this special solution and a free solution, such that A = 0. If boundary
conditions are imposed, the free solution may be quite difficult to find. In two dimensions, at
least, techniques of conformal transformation may be helpful.
Lower-dimensional cases
We consider the analogue of the previous problem in one and two dimensions,
starting with the latter. First, it must be made clear what is meant by twodimensional analogue. Here, it will be the study of a system described by a
potential field satisfying the dAlembert equation in dimension 2. Note that the
ensuing theory of electromagnetism is extremely different from the theory of
standard electromagnetism in dimension 3. Especially notable is that the
8
Disregarding the phenomenon of retinal persistence.

Note that clearly we should never write this integral in this manner, since distributions are
involved. We forgo absolute rigor in order to show a computation similar to those which are
exposed, often very quickly, in books of (classical and quantum) field theory.
9
419
Coulomb potential is not in 1/r but rather logarithmic,10 and the magnetic
field is not a vector but a simple scalar.11
What about the Green function of the dAlembertian? It is certainly possible to argue using a method identical with the preceding one, but it is simpler
to remark that it is possible to pass from the Green function in dimension 3
to the Green function in dimension 2 by integrating over one of the space
variables:
THEOREM 15.7 The Green functions of the dAlembertian in two and three dimen-
sions are related by the relation
G(2) (x, y, t) =
G(3) (x, y, z , t) dz ,
which may be seen as the convolution of G(3) with a linear Dirac source:
G(2) (x, y, t) = G(3) (x , y , z , t ) (x x) ( y y) (t t) 1(z ).
Proof. Recall that G(3) (x, y, z, t) = (x) ( y) (z) (t). Passing to the Fourier
transform (and putting c = 1 in the proof), we obtain the relation
(z 2 k x2 k 2y kz2 ) G (k x , k y , kz , z) = 1,
def
(15.7)
R +
which holds for all k x , k y , kz and z. We now put T (x, y, t) = G(3) (x, y, z , t) dz and
compute (2) T . Passing again to the Fourier transform, we obtain with T = F [T ]
that
(z 2 k x2 k 2y ) T (k x , k y , z) = (z 2 k x2 k 2y ) G(3) (k x , k y , 0, z),
which is equal to 1 according to equation (15.7) for kz = 0. Hence we get
(z 2 k x2 k 2y ) T (k x , k y , z) = 1,
or
(2) T (x, y, t) = (x) ( y) (t),
which is the equation defining the Green function of the dAlembertian in two dimensions, and therefore G(2) = T (up to a free solution).
We have therefore obtained

Z +
p

c
1
(ret)
G(2) (x, y, t) =
r 2 + z 2 c t dz ,
p
4
r 2 + z 2
where
10
def
r=
x 2 + y 2.
Since f ( r) = log k r k does satisfy f = 2 in two dimensions, see equation (7.11)

page 208.
11
See Chapter 17, page 483, concerning differential forms for an explanation of this seemingly bizarre phenomenon.
420
Green functions
Once more, this function is identically zero for negative values of t. For
positive values of t, we use the rule (explained in Theorem 7.45, page 199)

X 1
def

(z y)
f (z) =
with Z = y ; f ( y) = 0 ,
f ( y)
yZ
which is valid if the zeros of f are simple and isolated (i.e., if f ( y) does not
vanish at the same time as f ( y)). Applying this formula gives in particular
Z +
X g( y)

.
f (z ) g(z ) dz =
f ( y)
yZ
p
In the case considered we have f (z ) = r 2 + z 2 c t and
p
z
z1 = c 2 t 2 r 2 ,
f (z ) = p
Z = {z1 , z2 } with
z2 = z1 .
r 2 + z 2
After calculations, we find
c
p
2
2
2
2
(ret)
G(2)
(x, y, t) = 2 c t (x + y )
if c t > (x 2 + y 2 ),
(15.8)
otherwise.
A cylindrical outgoing wave can be seen which, in contrast to what happens

in three dimensions, is not a pulse. In figure 15.1, we show the radial function
(ret)
r 7 G(2)
(r, t) for two values t1 and t2 with t2 > t1 .
It should be remarked that, if we were living in a flat world,12 our visual
perception would have been entirely different. Instead of lighting a given
point (the eye of the subject) for a very short lapse of time, a light flash (for
instance, coming from a camera) would light it for an infinite amount of time,
p
although with an intensity rapidly decreasing13 with time (roughly like 1/ t,
as shown by formula (15.8)).
Remark 15.8 A number of electromagnetism books comment (or prove) the retarded potentials formula by noting that it is perfectly intuitive, the effect of a charge propagating at the
speed of light, etc. Unfortunately, this proof based on intuition is utterly wrong, since it completely disregards the phenomenon of nonlocality of the Green function of the dAlembertian
in dimension 2. It is true that the Green function in three dimensions is of pulse type, but it
is dangerous to found intuition on this fact!
The reader can check (using whichever method she prefers) that the Green
function of the dAlembertian in one dimension is given by
c /2 if c t > x,
G(1) (x, t) =
0
if c t < x.
12
As was envisioned by the pastor and academic Edwin Abbott (18391926) in his extraordinary novel Flatland [1].
13
Physiologically at least.
421
(ret)
G(2)
(r, t)
t1
t2
c t1
c t2
Fig. 15.1 The retarded Green function of the dAlembertian in two dimensions,
for
p
increasing values t2 > t1 of the variable t. In the abscissa, r = x 2 + y 2 .
15.3.c Covariant expression of advanced and retarded

Green functions
The advanced and retarded Green functions may be expressed in a covariant
manner. Denoting x = (c t, x) = (x0 , x) and x2 = x02 k xk2 = x02 r 2 , we
have (see, for instance, Exercise 8.9, page 242):

1
x2 = (c 2 t 2 r 2 ) =
(r c t) + (r + c t) ,
2r
which proves that

c
c
H (t) x2
and G (ad) ( x, t) =
H (t) x2 . (15.9)
G (ret) ( x, t) =
2
2
Although the function H (t) is not covariant, it should be noticed that
H (t) x2 is covariant (it is invariant under Lorentz transformations).14 This
is a fundamental property when trying to write the equations of electromagnetism or field theory in a covariant way [49].
15.3.d Radiation
What is the physical interpretation of the difference G (ad) G (ret) ?
We know that the vector potential is given, quite generally, by
ZZZZ
A(r) = free solution +
G (r r ) j(r ) d4 r ,
but there are two possible choices (at least) for the Green function, since either
the retarded or the advanced one may be used (or any linear combination
14
In fact, this expression is somewhat ambiguous since, for x = 0, the Dirac distribution
is singular and yet it is multiplied by a discontinuous function. The same ambiguity was,
however, also present in the first expression, since it was not specified what G was at t = 0...
Green functions
422
G (ad) + (1 )G (ret) ). The incoming and outgoing free solutions,15 denoted

A(in) and A(out) , respectively, are defined by the relations
ZZZZ
(in)
4
r R
G (ret) (r r ) j(r ) d4 r ,
(15.10)
A(r) = A (r) +
ZZZZ
r R4
G (ad) (r r ) j(r ) d4 r .
(15.11)
A(r) = A(out) (r) +
If the sources j are localized in space and time, then letting t in (15.10),
the integral in the right-hand side vanishes, since the Green function is retarded. The field A(in) can thus be interpreted as an incoming field, coming
from t = and evolving freely in time. One may also say that it is the
field from the past, asymptotically free, which we would like to see evolve since
t = without perturbing it.
Similarly, the field A(out) is the field which escapes, asymptotically freely
for t , that is, it is the field which, if left to evolve freely, becomes
(arbitrarily far in the future) equal to the total field for t +.
The difference between those two fields is therefore what is called the field
radiated by the current distribution j between t = and t = +. It is
equal to
ZZZZ
(ray)
(out)
(in)
A (r) = A
(r) A (r) =
G (r r ) j(r r ) d4 r ,
def
G (r) = G (ret) (r) G (ad) (r).
where
Hence the field A(ray) (r) gives, at any point and any time, the total field radiated by the distribution of charges. It is in particular interesting to notice
that, contrary to what is often stated, the unfortunate advanced Green function G (ad) is not entirely lacking in physical sense.
15.4
The heat equation
DEFINITION 15.9 The heat operator in dimension d is the differential oper-
ator
2
def
d
with d = 2 + + 2 ,
t
x1
xd
where c > 0 and > 0 are positive constants.
def
D=c
15
Recall that what is meant by free is a solution of the homogeneous Maxwell equations
(without source).
The heat equation
15.4.a
423
One-dimensional case
The Green function of the problem

We wish to solve the heat equation in one dimension, which is the equation

2
c 2 T (x, t) = (x, t),
(15.12)
t
x
where T (x, t) is the temperature depending on x R and t > 0, and where

is a heat source. The constant c > 0 represents a calorific capacity and > 0
is a thermal diffusion coefficient.
If no limit condition is imposed, the problem is translation invariant. We
will start by solving the analogue problem when the heat source is elementary
and modeled by a Dirac distribution. This means looking for the solution of
the equation

2
c 2 G (x, t) = (x, t),

t
x
and G will be the Green function of the heat equation which is also called
the heat kernel.16
Taking Fourier transforms, we have the correspondences
ZZ
G (x, t) 7 G (k, ) =
G (x, t) e 2i(t+k x) dt d,
7 2i
and
t
Hence the equation satisfied by G is
so that17
2
7 42 k 2 .
x2

42 k 2 + 2i c G (k, ) = 1,
G (k, ) =
(42 k 2
1
.
+ 2i c )
There only remains to take the inverse Fourier transforms. We start with the
variable . We have
Z
e 2it
Ge(k, t) =
d.
(42 k 2 + 2i c )
16
A kernel is a function K that can be used to express the solution f of a functional problem,
depending on a function g, in the form
Z
f (x) = K (x ; x ) g(x ) dx
(for instance, the Poisson, Dirichlet, and Fejr kernels already encountered). The Green functions are therefore also called kernels.

17
Up to addition of a distribution T satisfying 42 k 2 + 2i c T (k, ) = 0.
Green functions
424
If is considered as a complex variable, the integrand has a unique pole

located at 0 = 2ik 2 /c . Moreover, to evaluate this integral with the method
of residues, we must close the contour of integration of in the lower halfplane if t < 0 and in the upper half-plane if t > 0. The pole being situated in
the upper half-plane, it follows already that
G (x, t) 0
for t < 0.
This is a rather good thing, since it shows that the heat kernel is causal.
Remark 15.10 It is the presence of unique purely imaginary pole, with strictly positive imag-
inary part, which accounts for the fact that there is no advanced Green function here. In
other words, this is due to the fact that we have a partial derivative of order 1 with respect to t
and the equation has real coefficients.
In the case of quantum mechanics, the Schrdinger equation is formally identical to the
heat equation, but it has complex coefficients. The pole is then real, and both advanced and
retarded Green functions can be introduced.
For positive values of the time variable t, we close the contour in the upper
half-plane, apply Jordans second lemma, and thus derive

e 2it
2ik 2
e
G (k, t) = 2i Res
; =
(42 k 2 + 2i c )
c
=
e 4
2 k 2 t/c
(for t > 0).
Remembering that the Fourier transform of a gaussian is a gaussian, and

that the formula for scaling leads to
1 x 2 /a 2
2 2
F 1 e a k =
e
|a|
p
(which we apply with a = 2 t/c); the inverse Fourier transform in the
space variable is easily performed to obtain

1
c x 2
exp
G (x, t) = p
2 c t
4t
for t > 0.
(15.13)
Notice that, whereas temperature is zero everywhere for t < 0, the heat
source, although it is localized at x = 0 and t = 0, leads to strictly positive
temperatures everywhere for all t > 0. In other words, heat has propagated at
infinite speed. (Since this is certainly not literally the case, the heat equation
is not completely correct. However, given that the phenomena of thermal
diffusion are very slow, it may be considered that the microscopic processes
are indeed infinitely fast, and the heat equaton is perfectly satisfactory for
most thermal phenomena.)
The heat equation
425
Remark 15.11 The typical distance traveled by heat is x 2 = (4/c )t. This is also what
Einstein had postulated to explain the phenomenon of diffusion. Indeed, a particle under diffusion in a gas or liquid does not travel in a straight line, but rather follows a very roundabout
path, known as Brownian motion, which is characterized in particular by x 2 t. The
diffusion equation is in fact formally identical with the heat equation.
The gaussian we obtained in (15.13), with variance x 2 t, can also be interpreted
mathematically as a consequence of a probabilistic result, the central limit theorem, discussed
in Chapter 21.
The problem with an arbitraty heat source (15.12) is now easily solved, since
by linearity we can write
Z Z +
T (x, t) = [G ](x, t) =
G (x , t ) (x x , t t ) dt dx . (15.14)
Note, however, that this general solution again satisfies T (x, 0) = 0 for all
x R. We must still add to it a solution satisfying the initial conditions at
t = 0 and the free equation of heat propagation.
Initial conditions
We now want to take into account the initial conditions in the problem of
heat propagation. For this, we impose that the temperature at time t = 0 be
given by T (x, 0) = T0 (x) for all x R, where T0 is an arbitrary function.
We need a solution to the free problem (i.e., without heat source), which
means such that
T (x, t) = H (t) u(x, t).
The discontinuity of the solution at t = 0 provides a way to incorporate the
initial condition T (x, 0) = T0 (x). Indeed, in the sense of distributions we
have

T
u
=
+ T0 (x) (t).
t
t
The distribution T (x, t) then satisfies the equation

2
c
+ 2 T (x, t) = c T0 (x) (t),
t
x
the solution of which is given by
Z +
T (x, t) = c T0 (x) (t) G (x, t) = c
G (x x , t) T0(x ) dx .
Hence the solution of the problem that incorporates both initial condition
and heat source is given by
h
i
T (x, t) = T (x, t) + T (x, t) = (x, t) + c T0 (x) (t) G (x, t).
However, it is easier to take the initial conditions into account when using
the Laplace transform, as in the next example.
Green functions
426
15.4.b Three-dimensional case

We are now interested in the three-dimensional case, which we will treat using
the formalism of the Laplace transform.18
So we consider a function (or, if need be, a distribution) that satisfies the
free heat equation with a given boundary condition:
c
T
T =0
t
and
T ( r , 0) = T0 ( r ).
We perform a Laplace transform with respect to the time variable, denoted

T ( r , t) T ( r, p). We then have
(c p ) T ( r , p) = T0 ( r).
Notice that, this time, the initial conditions T0 ( r ) are automatically involved
in the Laplace transform, since the value of the temperature at t = 0 appears
in the differentiation theorem.
Taking the Fourier transform of this equation, we obtain
were we have denoted
f( k, p) = c Te ( k) Ge( k, p),
T
0
Ge( k, p) =
1
1
1
=
.
2
2
c p + 42 k 2 /c
c p + 4 k
(15.15)
A formal solution is then obtained by taking successively the inverse Laplace

transform:
Te ( k, t) = c Te0 ( k) Ge( k, t),
and then the inverse Fourier transform:
ZZZ

T ( r , t) = c T0 G ( r, t) =
c T0 ( r ) G ( r r , t) d3 r .
In order to obtain the function G ( r , t), we come back to the expression (15.15)
and use the table of Laplace transforms on page 614:
H (t) 42 k 2 t/c
e
Ge( k, t) =
c
whence (the inverse Fourier transform of a gaussian in three dimensions)

p
c
2
G ( r , t) =
e c r /4t .
(15.16)
8(t)3/2
18
The reader should not, however, get the impression that this is vastly preferable in three
dimensions: we could just as well use the method of the one-dimensional case, or we could
have treated the latter with the Laplace transform.
Quantum mechanics
427
Remark 15.12 We may add a right-hand side to the heat equation: D T ( r, t) = ( r, t), where
is the distribution describing the heat source; it allows us to derive a more general solution
(where the convolution is both on time and space variables)
h
i

T ( r, t) = + c T0 ( r) (t) G ( r, t).
Remark 15.13 Here also the typical distance of heat propagation is x 2 t. This result
(still a consequence of the central limit theorem!) is independent of the dimension of the
system.
15.5
Quantum mechanics
This time, we consider the evolution equation for a wave function:
DEFINITION 15.14 The Schrdinger operator is the differential operator
i h}
h}2
+
V ( r ),
t 2m
where m is the mass of the particle and V ( r ) is a potential independent of

time. If V 0, this is the free evolution equation.
In fact, only the free case V ( r ) 0 will be treated here. Notice that this
operator is formally similar to the heat operator with = i h}/2mc . However, its properties are very different, since in particular the heat equation
is irreversible, whereas the Schrdinger equation is perfectly reversible. This
is due to the fact that the Schrdinger equation is applied to complex-valued
functions.19
The Green function of the Schrdinger operator is computed in exactly
the same manner as in Section 15.420 and leads to
G ( p, t) = e i p
19
2 h} t/2m
= e i H t/ h}
for t > 0,
In the Schrdinger equation, it is therefore possible to perform the transformation t 7 t

and simultaneously to map 7 , so that the equation is still satisfied. Note that and have
the same physical interpretation (since the square of the modulus is the probability density of
the particle). Finally, notice that the change of into simply amounts to a change of
convention and to changing i into i, which cannot affect physics. The operator i h}/ t is
indeed invariant when simultaneously t 7 t and i 7 i.
20
This calculation involves a slight additional difficulty in that one must extract the square
root p
of a quantity which is this time imaginary. Hence a determination of the function
z 7 z must be specified, involving a cut in the plane. Once this is done, there is no further
difference. In the text, we have chosen a cut along the positive imaginary axis.
Green functions
428
(where we write H = p 2 /2m in momentum representation), taking p as

the conjugate variable to x and using the standard convention in quantum
mechanics:
ZZZ
G ( p, t) =
G ( x, t) e i p x/ h} d3 x.
Similarly, with the inverse Fourier transform, we obtain the equivalent
of (15.16):
G ( r , t) =
m 3/2
2
e im| r| /2 h}t
2i h}t
for t > 0.
(15.17)
If we want to solve the free Schrdinger equation

i h}
h}2
+
=0
t
2m
with initial conditions ( r , 0) = 0 ( r ), we get21 the relation for the propagation of the wave function:
Z
( r , t) = G ( r r , t) 0 ( r ) d3 r
for t > 0.
(15.18)
The function G is called the Schrdinger kernel. Equation (15.18) is therefore
perfectly equivalent to the Schrdinger equation written in the form
i h}
= H
t
def
with H =
h}2
.
2m
(15.19)
Either form may be used as a postulate for the purpose of the formalization
of quantum mechanics.
It should be noted that the general form of the Green function (15.17)
may be obtained through the very rich physical tool of path integration. The
reader is invited to look at the book of Feynman and Hibbs [34] to learn
more about this.
Note also that, using Diracs notation, we can write the evolution equation
of the particle as

d
i h} (t) = H (t) ,
dt
(denoting by | the wave vector of the particle), which can be solved
formally by putting

(t) = e i H t/ h} |0 .
21
Once more, those initial conditions are taken into accound by adding a right-hand side
given by 0 ( r ) (t).
Klein-Gordon equation
429
R
Using the closure relation 1 = | r r | d3 r (see Theorem 14.26) we obtain
for all t that
Z

( r , t) = r (t) = r e i H t/ h} r r 0 d3 r
Z
i H t/ h}
r 0 ( r ) d3 r .
=
r e
This is recognizable as the same integral as in (15.18). Since this relation

should be kept only for positive time, we are led to write

G ( r r , t) = H (t) r e i H t/ h} r .
In terms of Fourier transforms, we do indeed get

2
G ( p, t) = p e i H t/ h} p = e i p h}t/2m ,
which shows that the two expressions are consistent (and illustrates the simplicity of the Dirac notation).
15.6
In this section, the space and time variables will be expressed in a covariant
manner, that is, we will put x = (c t, x) = (x0 , x) and p = (p0 , p). We choose
a Minkowski metric with signature (+, , , ), which means that we put
p2 = p02 p 2 ,
x2 = x02 x 2 ,
and
2
.
t2
Remark 15.15 Note that the sign of the dAlembertian is opposite to that used in
Section 15.3; here we conform with relativistic customs, where usually = with
a metric (+, , , ).
We want to solve an equation of the type
( + m2 ) (x) = j(x),
which is called the Klein-Gordon equation for the field (it originates in
field theory). In this equation, m represents the mass of a particle, j is an
arbitrary source and is the field which is the unknown. As usual, we take
the constants h} and c to be equal to 1.22
22
Otherwise we should write ( h}2 c 2 + m 2 c 2 ) = j.
Green functions
430
To harmonize notation with those which are customary in field theory,

the Fourier transform and its inverse are denoted as follows:
Z
Z
1
i
px 4
e
f (p) = f (x) e d x and f (x) =
fe(p) e i px d4 p.
(2)4
We start by trying to solve the equation where the source term is replaced by
a Dirac distribution; the Green function of this equation
( + m2 ) G (x) = 4 (x)
is called a propagator. Taking the Fourier transform, we get
(p2 + m2 ) G (p) = 1,
or
G (p) =
p2 m2
1
.
p02 p 2 m2
As in the case of electromagnetism, there will arise a problem with this denominator when performing the inverse Fourier transform. Indeed, we will
try to write an expression like
ZZZZ
1
1
G (x) =
e i px d3 p d p0 .
2
4
(2)
p0 p 2 m 2
During
the process of integrating over the variable p0 , two poles at p0 =
p
p 2 + m2 occur. They must be bypassed by seeing p0 as a variable in the
complex plane and integrating along a deformed contour (see Section 15.2).
For instance, we may go around both poles in small half-circles in the
lower half-plane, which provides the advanced Green function:
G (ad) (t,
given by the path
x)
p 2 + m2
p 2 + m2
p0 C
This also amounts to adding a negative imaginary part to p0 , if integrating

first over p is preferred.
Similarly, the retarded Green function is defined by going around both
poles in the upper half-plane, or equivalently, by adding a small positive imaginary part to p0 :
G (ret) (t,
p0 C
given by the path
x)
Hence we put
def
G (ad) (p) =
(p0
p 2 + m2
1
p 2 m2
i)2
p 2 + m2
431
and
1
.
(p0
p 2 m2
Taking inverse Fourier transforms, we obtain
ZZZZ
1
e i px
(ad)/(ret)
G
(x) =
d4 p.
(2)4
(p0 i)2 p 2 m2
def
G (ret) (p) =
+ i)2
By the residue theorem, the function G (ret) is zero for t < 0 and G (ad) is
zero for t > 0. Moreover, these Green functions being defined by Lorentz
invariant objects, it follows that G (ret) is zero outside the future light cone (it
is said to be causal in the relativistic sense) and G (ret) outside of the past light
cone. Also, note that both functions are real-valued (despite the presence of
the term i) and that G (ret) (x) = G (ad) (x).
Whereas, in the case m2 = 0 of electromagnetism,23 we had a simple explicit formula for the Green functions, we would here get Bessel functions
after integrating on the variable p. The full calculation is not particularly
enlightening, and we leave it aside. Note, however, that the Green functions
are no longer supported on the light cone as with (15.9) (the case m2 = 0);
there is also propagation of signals at speeds less than the speed of light. This
is of course consistent with having a nonzero mass!
To close this very brief section (referring the reader to [49] for a wider
survey of the properties of the Green fnctions of the Klein-Gordon equation),
we mention that it is possible to introduce a very different Green function,
following Stueckelberg and Feynman. Let
GF (p) =
From the formula
x 2 2 + i
p2
1
.
m2 + i
x +
i
2

x +

i
,
2
which is valid to first order in , we deduce that changing m2 into m2 i

amounts to changing the poles
p
p
(1)
(2)
and
p0 = p 2 + m 2
p0 = p 2 + m 2
into, respectively,
p
p
p 2 + m2 + i
and
p 2 + m2 i
p
with = /2 p 2 + m2 , or to take the following integration contour:
23
Which is indeed a zero-mass field theory.
432
Green functions
given by the path
G F (t, x)
This last Green function is called the Feynman propagator24; it is complexvalued and does satisfy ( + m2 ) G F (x) = (x), as well as the relation
G F (x) = G F (x). Moreover, it can be shown (see [49]) that this propagator
is not zero outside the light cone (i.e., in particular that G F (0, r ) is not
identically zero for values of r 6= 0, contrary to what happens with the
advanced and retarded functions).
Remark 15.16 Why does this function not appear in classical physics? Simply because only
G (ad) and G (ret) are real-valued and have a physical meaning that is easy to grasp. The Feynman propagator G F is complex-valued; however, it makes sense in quantum mechanics, where
functions are intrinsically complex-valued.
EXERCISES
Exercise 15.1 Show that the Green function of the dAlembertian is of pulse type in odd
dimensions at least equal to 3, but that otherwise it is not localized.
Exercise 15.2 (Damped harmonic oscillator) A damped harmonic oscillator is ruled by the
following ordinary differential equation:

y + 2a y + 2 y = f ,
where a is some positive parameter. We are looking for the Green function associated to the
equation y + 2a y + 2 y = .
i) What is the equation which is satisfied by the Fourier transform of the Green function
Ge?
ii) Compute the Green function in D+ . One may represent graphically the poles of the
meromorphic function that appears during the calculations.
iii) Find the solution of the equation with an excitation f (some assumptions on f may
be required).
24
Richard Phillips Feynman (19181988), American physicist, was one of the inventors of
quantum electrodynamics (the quantum relativistic theory of electromagnetism), which is considered by many as one of the most beautiful current physical theories. With a penetrating
mind and a physical sense particularly acute and clear, he simplified the perturbative methods
and invented the diagrams that were named after him, which are extremely useful both for
making computations and for understanding them. His physics lecture notes are indispensable
reading and his autobiography is funny and illuminating in many respects.
Chapter
16
Tensors
In this chapter, we will give few details concerning the practical handling of tensors,
or the usual calculation tricks, juggling with indices, contractions, recognition of
tensors, etc. Many excellent books [5, 27, 35, 57, 62, 92] discuss those rules (and often
do much more than that!)
Here, the focus will be on the mathematical construction of tensors, in order to
complement the computational viewpoint of the physicist. This chapter is therefore not directed to physicists trying to learn how to tensor, but rather to those
who, having started to manipulate tensors, would like to understand what is hidden
behind.
Also, since this is a purely introductory chapter, we only discuss tensors in a flat
space. There will be no mention of parallel transport, fiber bundles, connections,
or Christoffel symbols. Readers wishing to learn more should consult any of the
classical textbooks of differential geometry (see Bibliography, page 617).
16.1
Tensors in affine space
Let K denote either the real field R or the complex field C. We identify
E = Kn , a vector space of dimension n, with an affine space E of the same
def
dimension, We consider E with a basis B = (e ) = (e1 , . . . , en ).
16.1.a
Vectors
Let u be a vector in E , that is, an element of Kn . There exists a unique

expression of u as linear combination of the basis vectors of B:
Tensors
434
DEFINITION 16.1 (Contravariant coordinates of a vector) The contravariant

coordinates of a vector u are the unique elements (u ) in K such that
n
P
u=
u e = u 1 e1 + + u n en .
=1
The contravariant coordinates are therefore the usual coordinates of linear

algebra.
Remark 16.2 Obviously, for a given vector u, the -th coordinate depends not only on the
vector e , but also on all other vectors of the chosen basis.
Remark 16.3 From the point of view of a physicist, what quantities may be modeled by a
vector? If working in R3 , a vector is not simply a triplet of real numbers; rather, it possesses
an additional property, related to changes of reference frame: if an observer performs a rotation,
the basis vectors turn also, and all triplets that are called vectors are transformed in a uniform
manner. A triplet (temperature, pressure, density), for instance, does not transform in the
same way, since each component is invariant under rotations (they are scalars).
In what manner are the coordinates of a vector transformed during a

change of basis? Let B = (e ) denote another basis of E . We write L the
matrix for the change of basis from B to B , that is, the matrix where the -th
column represents the vector e , as expressed in coordinates in the basis B
(see Appendix C). Writing this matrix L = (L ) , where is the line index
and the column index, it follows that for all [[1, n]] we have
e =
n
P
=1
L e .
(16.1)
The left index in L is the row index, and the right index is the column index.
For reasons which will become clearer later on, the first index is indicated in
superscript and the second in subscript. Symbolically, we write1
B = L B.
def
Let = L1 be the inverse matrix L, i.e., the unique matrix with coeffi
cients ( ) that satisfies
n
P
=1
L =
n
P
=1
L =
for all , [[1, n]]. The vectors of B can be expressed as functions of those
of B by means of the matrix :
1
Note that this is not a matrix relation. However, one can formally build a line vector with
the e s and write
(e1 , . . . , en ) = (e1 , . . . , en ) L.

B = B
or
e =
n
P
=1
435
e .
By writing the decomposition of a vector in the two bases, we can find the
transformation law for the coordinates of a vector:
P
P
u=
u e
= u e ,
u (
P
,
e )
u ( e ) =
P
,
P
,
u (L e ),
u (L e ),
(use (16.1))
(change names)
hence, by uniqueness of the decomposition in the basis B and B , respectively,

we have
u =
n
P
and
=1
u =
n
P
L u .
(16.2)
=1
Remark 16.4 These last formulas are also a criterion for being a vector: an n-uple, the com-
ponents of which are transformed according to (16.2) during a change of basis (for example,
when an observer performs a rotation) will be a vector. In what follows, saying that we want
to prove that a certain quantity is a vector will simply mean that we wish to establish the
relations (16.2).
16.1.b Einstein convention

Very quickly, the formulas for changes of coordinates become unreadable because of the increasing accumulation of indices of summation signs. This is
why Einstein2 suggested omitting the summation sign, which is implied each
time a Greek letter appears one time, and one time only, as a subscript and one time,
and one time only, as a superscript. Thus, relations (16.1) and (16.2) are expressed
as follows using the Einstein convention:
e = L e
and
e = e ,
u = u
and
u = L u .
(16.1 )
(16.2 )
Remark 16.5 One should be aware also that many books of relativity theory use the same letter
to denote both matrices and L. Two solutions are used to avoid conflicts:
In the book of S. Weinberg [92], for instance, they are distinguished by the position of
indices. Thus, represents what we denoted L , whereas is the same as ours.
2
Albert Einstein (18791955), German then Swiss, then American when, being Jewish, he
fled Nazi Germany physicist and mathematician, student of Hermann Minkowski, father of
special relativity (1905) then general relativity (around 1915), used and popularized the tensor
calculus, in particular by the introduction of riemannian geometry in physics.
Tensors
436
In the book of Ch. Misner, K. Thorne, and J. A. Wheeler [67], for instance, a primed
superscript indicates that is being understood, and a primed subscript that L is
understood.
Each notation has its advantages; the one we chose is closer to classical linear algebra. Here
is a short summary of various conventions in use, where is used to indicate any tensor object:
this book
Weinberg
Wheeler
= L
16.1.c Linear forms

We recall the following definitions:
DEFINITION 16.6 A linear form on E in any linear map from E to K.
DEFINITION 16.7 The dual of E , denoted E , is the vector space of linear
forms on E . For E , we denote either (v) or
, v the effect of
evaluating the linear form on a vector v E . Similarly, we sometimes write
, instead of .
Example 16.8 The coordinate linear forms
dx i :
Rn R,
(h 1 , . . . , h n ) 7 h i ,
defined for i = 1, . . . , n, are linear forms.
For a fixed vector u E , the map
: E K,
7
, u
is K-linear, hence it is an element in E . This application can only be

identically zero if u = 0. This shows that by associating it to u, we define an
injective map from E into E . In general, this is not a surjection. However,
it is so in two important cases: in Hilbert spaces (when considering only
continuous linear forms and linear forms defined by scalar products), and
when E is a finite-dimensional vector space:
THEOREM 16.9 (Bidual) Let E be a finite-dimensional K-vector space. Then we
have E E , dim E = dim E , and moreover E and E are canonically isomorphic.
THEOREM 16.10 (Dual basis) Let (e ) be a basis of E . The family (

) of
elements of E is defined by its action on any basis vector (e ) given by
, [[1, n]]
, e
1 if = ,
0 if 6= .
437
Then the family (

) is a basis of E , called the dual basis of (e ) I .
The linear forms (
) are also called coordinate forms, and often denoted
= dx ; to each vector u = u e , the linear form dx associates its coordinate u :

dx : u = u e 7 u .
Those two definitions are of course consistent.
Proof of Theorems 16.9 and 16.10. Let E be a linear form. For all [[1, n]],
denote = (e ). Then we have
(u) = (u e ) = u (e ) = (u),
for all u E , showing that = . Hence the family (
) is generating. It is
also free: if = 0, applying this form to the basis vector e yields = 0. So it is

a basis of E which consequently has the same dimension as E .
Since any linear form can be expressed uniquely in terms of the vectors of
the dual basis (
) , we make the following definition:
DEFINITION 16.11 (Covariant coordinates of a linear form) Let be a linear form, that is, an element in E . The covariant coordinates of the linear
form in the basis (
) are the coefficients appearing in its expression
in this basis:
=
(with the Einstein convention applied).

Hence the covariant coordinates of a linear form, just like the contravariant
coordinates of a vector, are its usual linear algebra coordinates.
They are computed using the following result:
THEOREM 16.12 (Expression of vectors and linear forms) Let (e ) be a basis
of E and let (
) be the dual basis of (e ) . Then, for any vector u E and any
linear form E , we have the decompositions

=
, e =
u=
, u e
u e
with =
, e
with u =
, u .
Proof. Indeed, we have
, u =
, u e = u
, e = u = u .
Similarly
, e = , e = = .
If the basis of E is changed, it is clear that the dual basis will also change.
Therefore, let e = L e , where L is an invertible matrix. We have then:
Tensors
438
PROPOSITION 16.13 If the basis of E is changed according to the matrix L, the

dual basis is changed according to the inverse matrix = L1 , that is, if we denote
(
) the dual basis of (e ) = L (e ) , we have
e = L e
and
for all .
Proof. We have indeed
, e = , L e = L = L = ,
which proves that the family ( ) is the dual basis of (e ) .
It is easy now to find the change of coordinates formula for a linear form
= , using the same technique as for vectors. We leave the two-line
computation as an exercise to the reader, and state the result:
PROPOSITION 16.14 The covariant coordinates of a linear form are changed in the
same manner as the basis vectors, namely,
= L .
Remark 16.15 The fact that the coordinates of a linear form are transformed in the same manner
as the vectors of the basis is the reason they are called covariant. Conversely, the coordinates of
a vector are changed in the inverse manner compared to the basis vectors (i.e., with the inverse
matrix ) and this explains why the coordinates are called contravariant coordinates.
Remark 16.16 Some quantities that may look like vectors actually transform like linear
forms. Thus, for a differentiable function f : E R, the quantity f transforms
as a linear form. The notation for the gradient is in fact a slightly awkward. It is
mathematically much more efficient to consider rather the linear form, depending on
the point x, given by
f
(x) dx
x
(this is an example of a differential form, see the next chapter). This linear form associates the
number d fx (h), which is equal to the scalar product f h, to any vector h E . Hence,
to define the vector f from the linear form d f |x , we must introduce this notion of scalar
product (what is called a metric). We emphasize this again: because a metric needs to be
introduced to define it, the gradient is not a vector.
d fx =
16.1.d Linear maps
Let : E E be a linear map. We denote ( ), its matrix in the basis B.

Note that the row index is indicated by a superscript, and the column index is indicated
by a subscript. The coefficients are given by

def
= , (e ) ,
Tensor product of vector spaces: tensors
439
and for any vector x, we have (x) = (x) e . Hence, under a change of
basis, the coordinates of of (i.e., its matrix coefficients) are transformed as
follows:
= , (e ) = , (L e ) = L .
PROPOSITION 16.17 (Change of coordinates) The coefficients of the matrix of
a linear map are transformed, with respect to the superscript, as the contravariant
coordinates of a vector, and with respect to the subscript, as the covariant coordinates of
a linear form.
Thus, if a map is representated by the matrix in the basis B and by in the
basis B , we have

= L .
This formula should be compared with the formula (c.2), on page 596.
Exercise 16.1 Show that the formulas
(x) = (x) e
and
(x) = (x) e
(Solution page 462).
are consistent.
16.1.e Lorentz transformations

As an illustration, we give here the matrices L and corresponding to the
Lorentz boost: the frame of (e ) has, relative to the frame of (e ), a velocity
v = v e x . Denoting = v/c and = (1 2 )1/2 , we have
0 0
0 0
0 0
0 0
.
L = (L ) =
= ( ) =
0
0
0
1 0
0 1 0
0
0
0 1
0
0 0 1
16.2
16.2.a
Existence of the tensor product of two vector spaces
Note: This section may be skipped on a first reading; it contains the statement
and proof of the existence theorem for the tensor space of two vector spaces.
THEOREM 16.18 Let E and F be two finite-dimensional vector spaces. There exists
a vector space E F such that, for any vector space G , the space of linear maps from
E F to G is isomorphic to the space of bilinear maps from E F to G , that is,
L (E F , G ) Bil(E F , G ).
Tensors
440
More precisely, there exists a bilinear map

: E F E F
such that, for any vector space G and any bilinear map f from E F to G , there
exists a unique linear map f from E F into G such that f = f , which is
summarized by the following diagram:
EF
@
@ f
@
@
?
R
@
- G
EF
The space E F is called the tensor product of E and F . It is only unique up to

(unique) isomorphism.
The bilinear map from E F to E F is denoted and also called the
tensor product. It is such that if (e i ) i and (f j ) j are basis of E and F , respectively,
the family (e i f j ) i, j is a basis of E F .
All this means that the tensor product linearizes what was bilinear, and also
that to any bilinear map on E F , one can associate a unique linear map on the
space E F .
Proof. See Appendix D, page 604.
It is important to remember that the tensor product of two vector spaces

is not defined uniquely, but always up to isomorphism. So there isnt a single
tensor product, but rather infinitely many, which are all isomorphic.
Exercise 16.2 Let n, p N. For A = (a i j ) Mn (C) and B M p (C), denote by A B the
matrix of size n p which is defined by blocks as follows:
a B a1n B
11.
..
AB =
..
. .
an1 B ann B
Show that the map (A , B) 7 A B gives an isomorphism
Mn (C) M p (C) M(n p) (C).
In particular, deduce that (Ei j Ekl )i, j,k,l is a basis of Mn (C) M p (C). Show that there exist
matrices in M(n p) (C) which are not of the form A B with A Mn (C) and B M p (C).
(Solution page 462).
THEOREM 16.19 If E and F are finite-dimensional vector spaces, we have
dim E F = dim E dim F .
441
PROPOSITION 16.20 In the particular case G = K, Theorem 16.18 gives
L (E F , K) Bil(E F , K).
For E and F finite-dimensional vector spaces, we have
E F Bil(E F , K),
that is, the tensor product E F is isomorphic to the space of bilinear forms on E F
(see Section 16.2.b).
16.2.b Tensor product of linear forms:

tensors of type 02
The abstract construction in the previous section is made more concrete when
considering the tensor product of E and F , where E and F are two finitedimensional vector spaces.
DEFINITION 16.21 Let E and F be two linear forms on E
and F , respectively. The tensor product of and , denoted , is the

bilinear form on E F defined by
: E F K,
(u, v) 7 (u) (v).
If (
) and (
) are bases of E and F , respectively, the family
(
)
is a basis of the vector space of bilinear forms on E F . In other words,

any element of Bil(E F , K) is a linear combination of the tensor products
. The space of bilinear forms on E F is a realization of the tensor
product E F defined in the previous section. This is denoted
Bil(E F , K) E F ,
and in general, with a slight abuse of notation, we identify those two spaces,
meaning that we do not hesitate to write an equal sign between the two.
Remark 16.22 Although it is perfectly possible to deal with the general case where E
are F are distinct spaces, we will henceforth assume that E = F , which is the most
common situation. Generalizing the formulas below causes no problem.

or 02 -tensor is any element in E
E , in other words any bilinear form on E E .
DEFINITION 16.23 A tensor of type
0
2
Tensors
442
Coordinates of a
THEOREM 16.24 Let T E E be a
the basis
with
of
by
0
-tensor
2
0
-tensor.
2
It can be expressed in terms of
T = T ,
def
T = T(e , e ).

Change of coordinates of a 02 -tensor
0
The coordinates of a 2 -tensor are transformed according to the following
theorem during a change of basis:
THEOREM 16.25 Let T E E be a tensor given by T = T
. During
a basis change, the coordinates T are transformed according to the rule
= L L T .
T

In other words, each covariant index of T is transformed using the matrix L.

Another viewpoint for a
0
-tensor
2
0
Since a 2 -tensor is (via the isomorphism in Theorem 16.20) nothing but a
bilinear form on E E , it is a machine that takes two vectors in E and gives
back a number:
T(vector, vector) = number.

Lets continue with this mechanical analogy. If we only put a single vector
in the machine, what happens? Then the machine expects another vector
(which we did not give) in order to produce a number. But a machine that
expects a vector to output a number is simply a linear form. Thus, once our
tensor is fed with one vector, it has been transformed into a linear form.
Precisely, if u E , the map
T(u, ) : v 7 T(u, v)
is indeed a linear form on E .
This means that the tensor T E E may also be seen as a machine
that turns a vector into a linear form:
T(, ) : E E ,
u 7 T(u, )
443
which does indeed take a vector (denoted u) as input and turns it to the
linear form T(u, ).
In mathematical terminology, showing that a bilinear form on E E
0
(a 2 -tensor), is the same thing as a linear map from E (the space of
vectors) to E (the space of linear forms) is denoted
Bil(E E , R) E E L (E , E ).
This property is clearly visible in the notation used for the coordinates of the
objects involved: if we write T = T , then we have
T(u, ) = T (u, )
, u
= T
= T u ,
(= T
, u
, )
which is a linear combination of the linear forms . The reader can check
carefully that those notations precisely describe the isomorphisms given above,
and that all this works with no effort whatsoever.
A particularly important instance of bilinear forms is the metric, and it
deserves its own section (see Section 16.3).
16.2.c Tensor product of vectors:

tensors of type 20
DEFINITION 16.26 Let u E and v E be two vectors. The tensor product
of u and v, denoted u v, is the bilinear form on E E defined by
u v : E E K,
(
, ) 7 (u) (v).
Any element in Bil(E E , K) is a linear combination of the family (e

e ) .
We denote E E = Bil(E E , K), and any element of E E is called

a tensor of type 20 or a 20 -tensor.
Coordinates of a
2
-tensor
0
The following theorem describes the coordinates of a
2
-tensor:
0
THEOREM 16.27 Let T E E be a tensor. Then T may be expressed in terms of

the basis (e e ) of E E as
T = T e e ,
with
def
T = T(
, ).
444
Tensors
Change of coordinates of a
2
-tensor
0
THEOREM 16.28 Let T E E be a tensor given by T = T e e . During
a basis change, the coordinates T are transformed according to the rule:

T = T .
In other words, each contravariant index of T is transformed using the matrix .
16.2.d Tensor product of a vector and a linear form:

linear maps or 11 -tensors
A special case of tensor product is the case of E E , since the previous
results allow us to write down isomorphisms
E E Bil(E E , K) L (E , E ) L (E , E ).
Therefore, it should be possible to see the tensor product of a vector and a

linear form as a linear map from E into itself!
DEFINITION 16.29 Let u E and E . The tensor product of the
vector u and the form is the bilinear form on E E defined by
u : E E K,
(
, v) 7
, u
, v .
Any bilinear form on E E may be expressed as linear combination of the

family (e ) . We write E E = Bil(E E , K) and any element of

E E is called a tensor of type 11 or 11 -tensor.
Now we can spell out concretely the statement above, namely that:
E E L (E , E ),
1
which gives an identification of a 1 -tensor and an endomorphism of E .
Let T = T e be such a tensor; then we can construct an associated

linear map from E to E by
v 7 T (v) e .
Conversely, a linear map may be seen as a
1
-tensor:
1
PROPOSITION 16.30 Let : E E be a linear map. Denote

def
= , (e )
445
A student of Jacobi and then of Dirichlet in Berlin, Leopold Kronecker (18231891) worked in finance for ten years starting when
he was 21. Thus enriched, he retired from business and dedicated
himself to mathematics. His interests ranged from Galois theory
(e.g., giving a simple proof that the general equation of degree at
least 5 cannot be solved with radicals), to elliptic functions, to
polynomial algebra. God created integers, all the rest is the work
of Man is his most famous quote. His constructive and finitist
viewpoints, opposed to those of Cantor, for instance, made him
pass for a reactionary. Ironically, with the advent of computers
and algorithmic questions, those ideas can now be seen as among
the most modern and lively.
its matrix coefficients in the basis B = (e ) . If (

) is the dual basis to B, then the
map can also be seen as a bilinear form on E E given by
= e .
Proof. Indeed, we can show that e is equal to by checking that it agrees

with when acting on a basis vector e . Since

e (e ) = , (e ) e (e ) = , (e ) e
, e

= , (e ) e = , (e ) e = (e ),
this is the case, and by linearity it follows that = e .

or 11 -tensor is any element of
E E or, equivalently, any linear map from E to itself.
1
1
Exercise 16.3 With the notation above, show that E E L (E , E ).
Notice, morevoer, that the coordinates of the identity map are the same in
any basis:
THEOREM 16.32 The coordinates of the identity map in any basis (e ) are given
by the Kronecker symbol
1 if = ,
def
=
0 if 6= .
Thus it is possible to write Id = and Id = e = e dx , independently

of the chosen basis.
Remark 16.33 The formula Id = e = e dx should
P be compared with the closure
relation given in Remark 9.26, page 260, namely Id =

|en en |. The tensor product is
implicit there, and the basis (en |) is the dual basis of (|en ) if the latter is orthonormal.
Tensors
446

Change of coordinates of a 11 -tensor
1
The coordinates of a 1 -tensor are neither more nor less than the coefficients of the matrix of the associated linear map. Therefore, Proposition 16.17
1
describes the law that governs the change of coordinates of a 1 -tensor.
THEOREM 16.34 Let T E E be a tensor given by T = T e . During
a basis change, the coordinates
are transformed according to the rule
T = L T .
Remark 16.35 In other words, each contravariant index of T is transformed using the matrix
and each covariant index is transformed using the matrix L. Compare this formula with the
change of basis formula (equation (c.2), page 596) in linear algebra.
16.2.e Tensors of type
p
q
Of course, it is possible to continue the tensor process an arbitrary number

of times, either with E as a factor or with E . One has to be careful with
the order of the factors, since the tensor product is not commutative, strictly
speaking. However, there are canonical isomorphisms E E E E (i.e.,
isomorphisms which are independent of the choice of a basis) which justify
some simplification, except that one must be quite careful when performing
index gymnastics, as described in the next section. It suffices therefore,
when dealing with higher-order tensors, to consider those of the type
p times
q times
z
}|
{ z
}|
{
E E E E E E ,
where the order of the factors is irrelevant.

def
DEFINITION 16.36 We denote E p = E E E , where the tensor
product is iterated p times, and similarly E q = E . . . E (q times).

p
p
or
-tensor, is an element of
q
q
p
q
the space E E
; in other words, it is a multilinear form on the space
E E E E.

Coordinates of a qp -tensor
p
THEOREM 16.38 Let T be a
-tensor. Then T may be expressed in terms of the
q
1
basis e1 e p q 1 p of the space E p E q :
T=
1 p
T1
q
1 q
e1 e p 1 q
The metric, or, how to raise and lower indices

with
1
T1
q
447

def
= T 1 , . . . , p , e1 , . . . , eq .
Change of coordinates of a
p
-tensor
q
1 p
THEOREM 16.39 Let T = T1 q e1 e p 1 q be a

p
1
tensor. During a basis change B = LB the coordinates T1
q
T 11q p =
For this reason, a
contravariant.
1
1
p
-tensor
q
p
p
1
1
q
q
p
q
are transformed
T11qp .
is said to be p times covariant and q times
16.3
The metric, or:
how to raise and lower indices
16.3.a
Metric and pseudo-metric
DEFINITION 16.40 A metric (or scalar product) on a finite-dimensional real
vector space E (K = R) is a symmetric, positive definite bilinear form; in other

words it is a bilinear map
g:
E 2 R,
(u, v) 7 g(u, v),
such that g(u, v) = g(v, u) for any u, v E , and such that g(u, u) > 0 for any
nonzero vector u.
A pseudo-metric on a finite-dimensional real vector space E is a symmetric
definite bilinear form (i.e., the only vector u for which we have g(u, v) = 0 for
all v E is the zero vector); however, g(u, u) itself may be of arbitrary sign
(and even zero, for a so-called light-like vector in special relativity).
Remark 16.41 A metric (resp. pseudo-metric) g can be seen, with the isomorphisms described
in the previous sections, as a
0
2
-tensor.
DEFINITION 16.42 (Coefficients g ) Let g be a metric on a real vector space

E of dimension n. Let B = (e1 , . . . , en ) be a basis of E and (
1 , . . . , n ) its
dual basis. The coordinates of g in the basis (

) of E E are
448
Tensors
Hermann Minkowski (18641909), German mathematician, was

born in Lithuania and spent his childhood in Russia. Since tsarist
laws prevented Jews from obtaining a decent education, his parents
emigrated to Germany. One of his brothers became a great physician, while he himself turned to mathematics. He taught in Bonn,
in Knigsberg, and then at the Zrich polytechnic school, where
Albert Einstein was his student. His mathematical works concern
number theory and mathematic physics. He gave a geometric interpretation of Einsteins special relativity, which is the most natural
framework, and his formalism, generalizing the notion of scalar
product and distance in space-time is still used today.
denoted g . In other words, we have
g = g dx dx .
Since g is symmetric, it follows that g = g for all , [[1, n]].
Example 16.43 The most important example of a metric is the euclidean metric which, in the
canonical basis (e ) of E , is given by coordinates
1 if = ,
g =
0 if 6= .
DEFINITION 16.44 A basis B = (e1 , . . . , en ) is orthonormal if
g(e , e ) =
for any , .
There can only be an orthonormal basis if the space is carrying a true metric,
and not a pseudo-metric (since the coordinates above describe the euclidean
metric in B).
DEFINITION 16.45 (Minkowski pseudo-metric) The Minkowski space is the
space R4 with the pseudo-metric defined by
if = = 0,
1
def
g = = 1 if = = 1, 2, or 3,
(16.3)
0
if 6= .
The space R4 with this pseudo-metric is denote M1,3 . This metric was introduced by H. Minkowski as the natural setting for the description of space-time
in einsteinian relativity.
DEFINITION 16.46 (Scalar product) The scalar product of the vectors u and
v in E is the quantity
The metric
449
def
u v = g(u, v).
If g is a pseudo-metric, we will also speak of pseudo-scalar product. It should

be noticed however that most books of relativity theory do not make an explicit
distinction; it is also customary to speak of minkowskian scalar product, or
more simply of minkowskian product.
PROPOSITION 16.47 The coordinates of the metric g are given by the scalar products
(defined by g) of the basis vectors
g = e e .
Proof. Indeed, we have e e = g(e , e ) = g (e , e )

= g , e
, e = g = g .
16.3.b Natural duality by means of the metric
Once a basis in E is chosen, we can exhibit an isomorphism between E and

E , which associates to a given basis vector e the dual vector , and which,
therefore, associates to an arbitrary vector v = v e the linear form
P
= v .
However, this isomorphism is not really interesting, because it depends on the

basis that has been chosen3 (the terminology used is that the isomorphism is
noncanonical.)
Using the metric tensor g, it is possible to introduce a duality of somewhat
different nature (a metric duality) between E and E which will be canonical
(i.e., it does not depend on the choice of a basis4).
u E on E
Indeed, if u E is a vector, we can introduce a linear form e
using the bilinear form g: it is given by
e
u : E R,
v 7 g(u, v) = u v
for any vector v. We denote this e

u = g(u, ).
DEFINITION 16.48 Let E be a real vector space with a metric g, and let B =
(e ) be a basis of E . The covariant coordinates of a vector u E are the
coordinates, in the basis (
) dual to B, of the linear form e
u E defined
by e
u, v = u v for all v E .
3
This artifical aspect is visible even in the way it is expressed: one

P has to abandon the
Einstein conventions and explicitly write a summation sign in = v , since the two
indices are at the same level. This is of course not a coincidence.
4
It depends on the metric, but the metric is usually given and is an intrinsic part of the
problem considered.
Tensors
450
u2
u2
e2
u 1 u1
e1
Fig. 16.1 Contravariant and covariant coordinates of a vector of R2 with the usual
euclidean scalar product. The vectors e1 and e2 have norm 1.
Those coordinates are denoted u , so that we have

e
u = u
whereas
u = u e .
THEOREM 16.49 Let u E be a vector. Its contravariant and covariant coordi-
nates are given respectively by
u =
, u = dx , u
and
u = u e = g(u, e ).
Proof. The first equality is none other than the definition 16.1 of the contravariant
coordinates of a vector. To prove
the second equality, remember that, by the very

definition of e
u, we have u e = e
u, e and therefore
u e = u , e = u , e = u = u .
It follows that the covariant and contravariant coordinates of a vector are

given by totally different means. Fix a vector u and choose a basis (e ) . The
-th contravariant coordinate of u in (e ), u = (u), is obtained using the
linear form , which depends not only on e , but on all vectors of the basis
(e ) . On the other hand, the covariant coordinate u is given by the scalar
product of u with e , which does not depend on the other basis vectors.
Remark 16.50 We see that the covariant coordinates correspond to orthogonal projections on the
axis spanned by the basis vector (requiring a metric to make precise the notion of orthogonality), whereas contravariant coordinates correspond to projections on the axis parallel to all
others (and the metric is not necessary). See figure 16.1.
Exercise 16.4 Show that, in an orthonormal basis, the covariant and contravariant coordinates are identical.
16.3.c Gymnastics: raising and lowering indices

THEOREM 16.51 Let u E be a vector with contravariant coordinates u and
covariant coordinates u . Then we have
u = g u .
The metric
451
Proof. Indeed, we find that u = u e = (u e ) e = u (e e ) = g u .
Hence the coefficients (g ), turn out to be useful to lower an index.

This does not apply only to vectors, but may be done for tensors of arbitrary
p
p1
type, always using the metric to transform a q -tensor into a q+1 -tensor.
However, beware! When performing such manipulations, the order of the
indices must be carefully taken care of (see Remark 16.60).
b E
THEOREM 16.52 Let E be a linear form. There exists a unique vector
such that
b v
, v =
for all v E .
This result comes from the definiteness of the metric and the fact that E
b
and E are isomorphic in finite dimensions. The explicit construction of
is described in Proposition 16.54.
DEFINITION 16.53 The contravariant coordinates of the linear are the
b . With a slight abuse of notation,

contravariant coordinates of the vector
they are denoted , so that we have
b = e .
PROPOSITION 16.54 Let E be a linear form. The following relations hold
between its covariant and contravariant coordinates:
b.
= , e
and
=
,
(16.4)
We have seen how the metric tensor can be used to lower the contravariant
index of a vector, thus obtaining its covariant coordinates. Similarly, it can be
used to lower the contravariant index of a linear form:
THEOREM 16.55 Let be a linear form. Its covariant and contravariant coordi-
nates satisfy the relations
= g .
b = e , and moreover for any vector v E , we have
Proof. By definition,

b v. In particular, = , e =
b e = e e = g , and this,
, v =
together with the relation g = g , concludes the proof.
We now pass to the question of raising indices.
DEFINITION 16.56 (Coefficients g ) Let g = g dx dx be the metric
tensor. Its dual

is defined by
def
b
g = g e e
b
b , b ) =
b b
g(
, ) = g(
, E .
Tensors
452
In particular, the coefficients g satisfy g = b

g(
, ) by Theorem 16.27.
2
Moreover, we have b
g E E ; that is, the dual of the metric is a 0 -tensor.
The dual first appears in the coordinates of the basis vectors of E and E ,
which are given by the following proposition:
PROPOSITION 16.57 The covariant and contravariant coordinates of e and
denoted [e ] , with square brackets are given by
[e ] = g ,
] = g ,
[
[e ] = ,
[
] = .
Proof. The first equality comes directly from v = v e applied to v = e . The

second equality holds because
(
) is the -th contravariant coordinate of the vector
c , and is therefore equal to , c = c c = g . The third equality is trivial

since the -th coordinate of e is , e . Finally, the last equality follows from

(
) = , e .
THEOREM 16.58 The matrices g and g are inverse to each other, that is, we
have
g g =
and
g g = .
Proof. First we prove the second formula. According to Proposition 16.57, we have
g = (
) , hence g g = (
) g = (
) by Theorem 16.55. Then the desired
conclusion follows from Proposition 16.57 again.
The first equality is implied by the second.
From this we can see that the dual tensor b

g = g e e can be used to
raise indices:
THEOREM 16.59 Let u E be a vector and let E be a linear form. Then
their covariant and contravariant coordinates satisfy the following relations:

u = g u
and
= g .
Proof. This is clear from to the formulas related to lowering indices, since g and
g are inverse to each other.
Of course these results are also valid (with identical proofs!) for
p+1
tensors, which are transformed into q1 -tensors by b
g (i.e., by g ).
p
q
Remark 16.60 This is where the abuse of notation consisting in putting together all
covariant indices (and all contravariant indices) becomes dangerous. Indeed, consider,
for instance, a tensor with coordinates
T ,
Operations on tensors
453
in expanded notation, and suppose now that, to simplify notation, we denote the coordinates
simply
T ,
using the fact that E E E is canonically isomorphic to E 2 E .

Now, suppose we want to lower the index by contracting it with the metric tensor.
Should the result be denoted
T or T
?
Only the original noncondensed form can give the answer. If the metric tensor is used, it is
important to be careful to write all indices properly in the correct position.
16.4
Operations on tensors
Theorem 16.59 may be generalized, noting that the tensors g and b
g are
p
p+1
p1
used to transform q -tensors into q1 -tensors or q+1 -tensors, respectively.
For instance, suppose we have a tensor denoted T in physics. We may
lower the index using g , defining
def
T =g T .
This is the viewpoint of index-manipulations. What is the intrinsic mathematical meaning behind this operation?
The tensor T, written
T = T e e ,
is the object which associates the real number
T(
, , v) = T v
to a triple (
, , v) E E E . As to the tensor where one index is
lowered, it corresponds intrinsically to
e = T e : (
T
, u, v) 7 T u v .

e , it suffices simply to be able to pass from u to . But this is

To define T
precisely the role of g, since from u E we may define a linear form e
u E
by the relation
e
u : E R,
w 7 g(u, w).
e as the object which

Then we can define T
1. takes as input the triple (
, u, v),
(16.5)
Tensors
454
2. transforms it into (
, e
u, v), where e
u is defined by the relation (16.5),
u, v):
3. outputs the real number T(
, e
e : E E E R,
T
(
, u, v) 7 T(
, e
u, v).
The same operation may be performed to raise an index. This time, the
tensor b
g is used in order to define a vector, starting with a linear form (still
because of the metric duality).
To summarize:
THEOREM 16.61 (Raising and lowering indices) The metric tensors g and b
g can
be used to transform a
p
-tensor
q
into a
p+1
-tensor
q1
or a
p1
-tensor,
q+1
respectively.
There is another fundamental operation defined on tensors, namely, the

contraction of indices. Once again, to start with a concrete example, consider
a tensor
T = T e e .
2
1
We will associate to this 1 -tensor a 0 -tensor, by defining
def
T = T .
Thus we have performed an index contraction. For instance, the trace of a

linear operator in a finite-dimensional vector space is given by the contraction
of its indices.
Intrinsically, we are looking for a new tensor
e = T e : 7 T .
T
) , we have
Notice that, taking an arbitrary basis (e ) and its dual basis (
T(
, , e ) = T [
] [e ] = T = T
= T .
e by
Hence it suffices to define the tensor T
e : E R,
T
7 T(
, , e ),
e = T e with T = T holds indeed (indepenand then the relation T
dently of the choice of a basis).
THEOREM 16.62 (Index contraction) There exists an operation of index contrac-
tion which, for any choice of one covariant and one contravariant coordinate, transp
p1
forms a q -tensor into a q1 -tensor. This operation does not require the use of the
metric, and is defined independently of the choice of a basis.
Change of coordinates
455
16.5
One of the most natural question that a physicist may ask is how does a
change of coordinates transform a quantity q?
First, notice that a change of coordinates may be different things:
a simple isometry of the ambient space (rotation, translation, symmetry,...);
a change of Galilean reference frame (which is then an isometry in
Minkowski space-time with its pseudo-metric);
a more complex abstract transformation (going from cartesian to polar
coordinates, arbitrary coordinates in general relativity,...).
16.5.a
Curvilinear coordinates
Consider general coordinates (x, y, z) in R3 , which we rather denote (x 1 , x 2 , x 3 ),

or even X = (x 1 , . . . , x n ) in order to easily generalize to a space with n dimensions. For instance, this may be euclidean coordinates in Rn .
Suppose given a set of functions which define new coordinates, called
curvilinear coordinates:
1
u = u 1 (x 1 , . . . , x n ),
.
..
.
U=
.
.
u n = u n (x 1 , . . . , x n ).
These functions define a C 1 change of coordinates if (and only if) the map
: Rn Rn
(x 1 , . . . , x n ) 7 (u 1 , . . . , u n )
is of C 1 -class (its partial derivatives exist and are continuous), it is injective

and its image is open, and if 1 is of C 1 -class on this open subset (this is
called a C 1 -diffeomorphism in mathematics, but physicists simply speak of
coordinate change). Recall the following result of differential calculus [75]:
is a C 1 -diffeomorphism if and only if it is of C 1 -class and its jacobian is
nowhere zero. In other words, its differential d is a bijective linear map at
every point.
How is a vector such as a velocity vector, for instance, transformed? Its
norm, expressed in kilometers per hour (hence, in terms of a first set of
coordinates) is not the same expressed in miles per second (in a second set
of coordinates). However, the velocity is a physical quantity independent of
the manner in which it is measured. Therefore, the intrinsic physical quantity
must not be mistaken for its coordinates in a given system of measurement.
Tensors
456
Unfortunately, physicists often use the same word for those two very different
concepts! Or, which amounts to the same thing here, they omit to state
whether active or passive transformations are involved.
In what follows, we consider passive transformations: the space is invariant,
as well as the quantities measured; on the other hand, the coordinate system
varies, and hence the coordinates of the quantities measured also change.
16.5.b Basis vectors

Let Y be a point in R, and let e1 (Y), . . . , en (Y) denote the basis vectors
in Rn attached to the point Y. This family does not depend on Y, since X
represents the canonical coordinates on Rn .
x2 = 2
x2 = 1
x2 = 3
x1 = 3
x1 = 2
e2
Y
x1 = 1
e1
x1 = 0
We now want to find the new vectors, attached at the point Y = (Y) and tan-
gent to the coordinate lines, which represent the vectors e1 (Y), . . . , en (Y)
after the change of coordinates. They will then be the new basis vectors. We
denote them e1 (Y ), . . . , en (Y ) .
u2 = 3
u2 = 2
u2 = 1
u2 = 4
u1 = 2
u1 = 1
e1
e2
e2
Y
e1
u1 = 0
457
How are they defined? Suppose that we have two nearby points M and
M in space. Then we can write
M M = xm e ,
where xm is the variation of the -th coordinate x when going from M
to M . This same vector, which is an intrinsic quantity, may also be written
M M = um e ,
where um is the variation of the -th curvilinear coordinate u when going

from M to M . But, to first order, we have the relation
um =
u
xm ,
x
which, after identifying, gives
xm e =
u
u
x
=
xm e
m
x
x
(since and are mute indices), and hence, since the xm are independent
and arbitrary
u
x
e
or conversely
e =
e ,
x
u
or again, expanding the notation,
1

un
u
e1
e
1
1
1
x
x
.. ..
..
.
.
=
. .
. ,
.
1

u
un
en
en
xn
xn
and
1

x
xn
e1
e
1
1
1
u
u
..
.. ..
..
. = .
. .
.
1

n
x
en
en
un
un
Note that the matrices
u
x
=
and
L =
x
u
e =
represent, respectively, the jacobian matrices of and 1 . In particular, they

are invertible (the differential of is bijective and therefore invertible), and
each is in fact the inverse of the other. This may be written as
Tensors
458
u x
= .
x u
We recover the notation used in an affine space5; the only difference is that,
now, the matrices L and depend on the point M where the transformation
is performed. To summarize:
THEOREM 16.63 (Transformation of basis vectors) During a change of coordi-
nates : (x ) 7 (u ) , the basis vectors are transformed according to the following

relations:
x
u
e =
e
and
e =
e .
u
x
Note that these formulas are balanced from the point of view of indices,
obeying the following rule: a superscripted index of a term in the denominator of a
fraction is equivalent to a subscripted index.
16.5.c Transformation of physical quantities

Now, assume we have a vector which corresponds to some physical quantity
(and hence is an intrinsic object, also called a vector field). It can be expressed
by
v = v e = v e
in each system of coordinates. Applying the rule for the transformation of
basis vectors, we derive the rule for the, transformation of the contravariant
coordinates of a vector:
THEOREM 16.64 (Transformation of contravariant coordinates)
During a change of coordinates : (x ) 7 (u ) , the contravariant coordinates v

of a vector field are transformed according to the relation
v =
u
v
x
and
v =
x
v ,
u
that is, following the opposite rule as that used for basis vectors.6
5
Consider in this manner the formula of Proposition 16.13. We may write = in

the form dx = dx , which gives the well-known formula
dx =
x
dx .
x
6
One should of course indicate every time at which point the vectors are tangent, which is
also where the partial derivatives are evaluated. In this case, the preceding formulas should be
459
Remark 16.65 Here is a way to memorize those formulas. Let x be the new coordinates. We
recover the relation expressing e as a function of e , namely,

e =
x
e,
x
by applying the following rules:

since the index of the vector on the left-hand side of the equation (which is e ) is
subscripted, we write e = xx e and put the prime sign at the bottom:
x
e;
x
e =
since this index is subscripted, it can appear as x only in the denominator, so we

write
e =
x
e;
x
finally, we complete the indices with a mute :

e =
x
e
x
(and then one should check that the index is balanced as it should).
Similarly, to recover the formula giving v in terms of the v s,
we write the formula with a prime sign in the numerator, since the index is superscripted:
v =
x
v ;
x
then we complete the indices in order that the result be balanced:

v =
x
v .
x
The reader is invited to check that the converse formulas (giving v and e in terms of the
and e s) may be obtained with similar techniques.
v s
16.5.d Transformation of linear forms

Since linear forms are dual to vectors, the reader will not find the following
result very surprising, except for the terminology: what might be called a field
of linear forms is called a differential form:
written more rigorously:
u
(Y) v (Y)
and
x
which complicates the notation quite a bit.
v (Y ) =
v (Y) =
x
(Y ) v (Y ),
u
Tensors
460
THEOREM 16.66 (Transformation of linear forms) Let be a differential form.
During a change of coordinates : (x ) 7 (u ) , its covariant coordinates are

transformed according to the relations
=
and
u
.
x
x
.
u
The basic linear forms are transformed according to

=
and
The rule sketched above to memorize those relations remains valid.
16.5.e Transformation of an arbitrary tensor field

Of course, the previous results extend easily to tensors of arbitrary type.
THEOREM 16.67 (Transformation of a tensor) Let T be a tensor field with coor-
dinates
...,
1
p
T(Y) = T1 ,...,
e1 (Y ) en (Y ) 1 (Y) q (Y).
q
During a change of coordinates : (x ) 7 (u ) , these coordinates are transformed

...,
T 11,...,qp =
u 1
u p x 1
x q 1 ,..., p
T1 ,...,q .
x 1
x p u 1
u q
Finally, we note that it is customary, both in relativity theory and in field

theory, to introduce the following notation:
def
and
def
T, = T ,
so that, for instance, the divergence of a vector T = T e is written

div T = T , .
461
16.5.f Conclusion
Le Calcul Tensoriel sait mieux la physique que le Physicien lui-mme.
Tensor calculus knows physics better than the physicist himself does.
Paul Langevin (quoted in [10]).
Physical applications are too numerous to be studied individually. Among

others, we may mention:
Special relativity: the natural setting of special relativity was defined by Minkowski. It is a space isomorphic to R4 , with a pseudo-metric with
signature (1, 3). Generalizing the scalar product is done very simply.
The notion of covariance, in other words, of the invariance of a quantity
or a the form of an equation, is expressed mathematically in the invariance under change of basis of quantities which are balanced in terms
of their expressions using indices.
General relativity: Riemannian geometry is unfortunately outside the scope
of this book. We simply mention that the notion of Minkowski space
is extended to non-euclidean geometry. Tensor calculus (in the context
of manifolds) is then, for intrinsic reasons, an indispensable ingredient
of the theory.
Fluid mechanics: using tensors, the notion of current can be generalized.
The current of a vector is a scalar. The current of a vector may be
expressed in matrix form (one index for each component of the vector,
one index for the direction of the current). It is then convenient to
go farther and to denote as a tensor with three indices the current of a
tensor with two indices, and so on.
Electricity, signal theory, optics: tensor calculus may be used as a mathematical trick to express general linear transformations. Then the underlying
notions such as linear forms, metric, duality,... do not carry special physical significance.
Remark 16.68 A reader interested in general tensor calculus is invited to read the books of
Delachet [27]; the excellent treatise on modern geometry [30] in three volumes; the bible
concerning the mathematics of general relativity (including all non-euclidean geometry) by
Wheeler7 et al. [67]; or finally, one of the best physics books ever, that of Steven Weinberg [92].
On voit [...] se substituer lhomo faber lhomo mathematicus.
Par exemple loutil tensoriel est un merveilleux oprateur de gnralit ;
le manier, lesprit acquiert des capacits nouvelles de gnralisation.
We see [...] homo mathematicus substituting himself to homo faber.
For instance, tensor calculus is a wonderful operator in generalities; handling it,
the mind develops new capacities for generalization.
Gaston Bachelard, Le nouvel esprit scientifique [10]
7
Who invented the words black hole.
Tensors
462
SOLUTIONS
Solution of exercise 16.1 on page 439. Using the formulas for change of basis, we get

(x) e = L (x) L e

= L L (x) e
= (x) e = (x) e ,
and this is the required formula.
Solution of exercise 16.2 on page 440. The first questions are only a matter of writing
things down. In the case n = p = 2, we see, for instance, that
0 1
0
1
2 3 2 3 1 1 0
0 0
0
2
0
2
2
0 0
4
6
but
0
0
0
0
cannot be a tensor product of two matrices.
0
1
0
0
0
0
2
0
0
0
3
7
1
,
3
Chapter
17
Differential forms
In all this chapter, E denotes a real vector space of finite dimension n (hence
we have E Rn , but E does not, a priori, have a euclidean structure, for
instance).
17.1
Exterior algebra
17.1.a 1-forms
DEFINITION 17.1 A linear form or exterior 1-form on E is a linear map
from E to R:
: E R.
The dual of E is the vector space E of linear forms on E .
Let B = ( b 1 , . . . , b n ) be a basis of E . Any vector x E can be expressed
in a unique way as combination of vectors of the basis B:
x = x 1 b1 + + x n bn.
We often denote simply x = (x 1 , . . . , x n ) (identifying E and Rn , the basis

of E being assumed given). For any i [[1, n]], we denote either b i or dx i
the linear form on E which is defined by
i

b i ( x) = b i , x = x i
or
b , b j = ij
i, j [[1, n]].
Differential forms
464
THEOREM 17.2 The family (b 1 , . . . , b n ) is a basis of E .
Proof.
i ), we have obviously (by linearity)
P Let E . Then, denoting i = (b
= i i b i , which P
shows that the family (b 1 , . . . , b n ) generates E . It is also a
free family (indeed, if i b i = 0, applying this linear form to the vector b j yields
j = 0), and hence is a basis.
Any linear form may therefore be written

= 1 dx 1 + + n dx n
with i = ( b i ).
17.1.b Exterior 2-forms

DEFINITION 17.3 An exterior form of degree 2 or a 2-form is an antisymmetric bilinear form on E E , that is, a bilinear form 2 : E E R such
that we have
for all x, y E .
2 ( x, y) = 2 ( y, x)
Example 17.4 The determinant map in a basis B = ( b 1 , b 2 ) of R2 , namely,
det : E E R,

x 1
( x, y) 7 det( x, y) = 2
x
B
is an exterior 2-form on R2 .

y 1
,
y2
DEFINITION 17.5 The vector space of exterior forms of degree 2 on E is
denoted 2 (E).
We now look for a basis of 2 (E ). Clearly, the bilinear form

dx 1 dx 2 dx 2 dx 1 : E E R,
( x, y) 7 x 1 y 2 x 2 y 1
is an element of 2 (E ). Similarly, the bilinear forms dx i dx j dx j dx i

with i, j [[1, n]] are also in this space. A moment of thought suffices to
see that any exterior 2-form may be expressed as linear combination of those
elements. So we have found a generating family. But clearly, this family is not
free, since
dx i dx i dx i dx i = 0 for all i [[1, n]];

dx j dx i dx i dx j = dx i dx j dx j dx i for all i, j [[1, n]].
It is sufficient to keep the elements of this type with i < j; there are n(n1)/2
of them.
THEOREM 17.6 A basis of 2 (E ) is given by the family of the exterior 2-forms
def
dx i dx j = dx i dx j dx j dx i ,
where 1 i < j n. Hence we have dim 2 (E ) = n(n 1)/2.
Exterior algebra
465
def
Notice that the space E 2 = E E admits for basis the family of

elements
i

dx dx j ; i, j [[1, n]] ,
from which it is natural to create two further families:
the family we have seen of elements of the type dx i dx j (with i < j);
the family of elements of the type 12 (dx i dx j + dx j dx i ).
Those two families are bases, respectively, for the space 2 (E ), which is the
0
0
space of antisymmetric 2 -tensors, and for the space of symmetric 2 -tensors.
PROPOSITION 17.7 Any
tensor.
This may be generalized to
0
-tensor
2
is the sum of a symmetric and an antisymmetric
0
-tensors.
k
17.1.c Exterior k-forms

We first recall the following definition:
DEFINITION 17.8 A permutation of [[1, n]] is any bijection of the set [[1, n]]
into itself; in other words, permutes the integers from 1 to n. A transposition is any permutation which exchanges two elements exactly and leaves the
others invariants. The set of permutations of [[1, n]] is denoted Sn .
The composition of applications defines the product of permutations.
Any permutation may be written as a product of finitely many transpositions. The signature of a permutation Sn is the integer () given by:
() = 1 if can be written as a product a an even number of transpositions,
and () = 1 if can be written as the product of an odd number of
transpositions.1 For any permutations and , we have
() = () ().
DEFINITION 17.9 An exterior k-form, or more simply a k-form, is any an-
tisymmetric k-multilinear form k on E k , that is, the map k : E k R

satisfies the relation
k ( x 1 , . . . , x k ) = (1)() k ( x (1) , . . . , x (k) )
for any permutation Sk .
1
Implicit in this definition is the fact, which we admit, that the parity of the number
of transpositions is independent of the choice of a formula expressing as a product of
transpositions.
Differential forms
466
Example 17.10 The most famous example is that of an exterior k-form on a space of dimension
k, namely, the determinant map expressed using the canonical basis: if ( e 1 , . . . , e k ) is the
canonical basis of Rk , and if i Rk for i = 1, . . . , k, then denoting i = i1 e 1 + + ik e k ,
this is a map
det :
(Rk )k R,

11

( 1 , . . . , k ) 7 detcan ( 1 , . . . , k ) = ...

k1
which is an exterior k-form.
...
...

1k
.. ,
.
kk
Example 17.11 For a fixed vector u R3 , the map : ( x, y) 7 det( x, y, u) on R3 R3 is
an exterior 2-form.
DEFINITION 17.12 The vector space of exterior k-forms on E is denoted
k (E).
We will now find a basis of k (E ).
To generalize the corresponding result of the preceding section for 2 (E ),
we introduce the following definition:
DEFINITION 17.13 Let k N and let i1 , . . . , ik [[1, n]]. Then the exterior
k-form dx i1 dx ik is defined by the formula
X
def
dx i1 dx ik =
() dx (i1 ) dx (ik ) .
Sk
Example 17.14 Consider the space R4 and denote by x, y, z, t the coordonnates in this space.
Then we have
dx dy dz = dx dy dz dx dz dy + dy dz dx dy dx dz
+ dz dx dy dz dy dx.
THEOREM 17.15 The family
i

dx 1 dx ik ; 1 i1 < < ik n
n
is a basis of k (Rn ), the dimension of which is equal to the binomial coefficient k .
0
Hence any exterior k-form (i.e., any antisymmetric k -tensor) can be written in the
form
P
k =
i1 ,...,ik dx i1 dx ik .
i1 <<ik
In the remainder of this chapter, Einsteins summation conventions will

be extended to exterior forms, so we will write
P
= i1 ,...,ik dx i1 dx ik instead of = i1 ,...,ik dx i1 dx ik ,
i1 <<ik
Exterior algebra
467
that is, when no ambiguity arises, repeated indices (appearing once as superscript and once as subscript) are assumed to be summed
in strictly increasing
n
order; so if there are k indices, the sum is over k terms instead of nk terms.
For the special case k = n, we deduce the following result:
THEOREM 17.16 The space n (Rn ) is of dimension 1.
COROLLARY 17.16.1 Any exterior n-form on Rn is a multiple of the form deter-
minant in the canonical basis.
DEFINITION 17.17 A volume form on Rn is any nonzero exterior n-form.

Remark 17.18 We will often denote k an exterior k-form, p an exterior p-form,
and so on. This supscript notation is only intended to recall the degree of the form
which is considered. In particular, it should not be mistaken with the exponent in a
power expression.
17.1.d Exterior product

There is a product defined on the various spaces of exterior forms, which
associates a (k+)-exterior form k to a pair (k , ) in k (E ) (E ).
We will first give the definition for the product of two 1-forms, then generalize
it. Note that the notation will be consistent with the previous introduction
of basis elements dx i1 dx ik for k (E ) (i.e., those will be products of
1-forms).
Exterior product of two 1-forms
DEFINITION 17.19 Let , be exterior 1-forms. The exterior product
of and is defined by
= ( ).
def
It is easily checked that is an exterior 2-form.
PROPOSITION 17.20 Let and be exterior 1-forms. Then we have
)( x,

( x) ( x)

,
y) =
( y) ( y)
and moreover = . The exterior product is bilinear and antisymmetric.

Similarly, the product of p exterior 1-forms may be defined as follows:
DEFINITION 17.21 Let 1 , . . . , p be exterior 1-forms. Then the exterior
product 1 p p (E ) is defined as the multilinear map on E p
Differential forms
468
such that

1 ( x 1 ) . . . p ( x 1 )

def
.. .
1 p ( x 1 , . . . , x p ) = ...

.

1 ( x p ) . . . p ( x p )
Exterior product of two arbitrary exterior forms
We now introduce the general definition of the exterior product of two exterior forms.
DEFINITION 17.22 (Exterior product) Let k, N , k k (E ) and
(E ). The exterior product of k and is the exterior (k + )-form

denoted k such that
X
def
k ( x 1 , . . . , x k+ ) =
() k ( x (1) , . . . , x (k) )
Sk+
( x (k+1) , . . . , x (k+) ).
Example 17.23 For a 2-form 2 and a 1-form 1 , we have
2 1 (a, b , c ) = 22 (a, b ) 1 (c ) + 22 (b , c ) 1 (a) + 22 (c , a) 1 (b ).
The reader is invited, as an exercise, to check the following result:

PROPOSITION 17.24 The exterior product is
(anti)symmetric: p q = (1) p q q p for p p (E ) and

q q (E );
distributive with respect to addition: ( + ) = + ;

associative: ( ) = ( ).
Example 17.25 Let and be exterior 1-forms. Then using Definition 17.22 we get
( x, y) = ( x) ( y) ( x) ( y),
i.e. = as defined in the preceding paragraph.
PROPOSITION 17.26 Let 1 , . . . , p , p+1 , . . . , p+q be exterior 1-forms. Then
we have
(1 p ) ( p+1 p+q ) = 1 p p+1 p+q .

Proof. This formula follows from the definition of the exterior product of 1-forms
and the compatibility with the general definition.
Differential forms on a vector space
469
DEFINITION 17.27 The exterior algebra on Rn is the vector space
(Rn ) =
def
n
M
k (Rn ).
k=0
The exterior product is extended to the exterior algebra by distributivity.
17.2
Differential forms on a vector space
17.2.a
Definition
As before, in all this section, E = Rn is a real vector space of finite dimension n.

DEFINITION 17.28 Let p be an integer or +.
A differential form of
degree k on R n of C p -class, or differential k-form, is a map of C p -class
: Rn k (Rn ),
x 7 ( x).
For any point x in Rn , ( x) is therefore an exterior k-form on Rn .

What is the precise meaning of C p -class here? Any differential k-form
may be written
( x) = i1 ,...,ik ( x) dx i1 dx ik ,
where x 7 i1 ,...,ik ( x) are real-valued functions on Rn . Then x 7 ( x) is

n
of C p -class if and only if the k functions x 7 i1 ,...,ik ( x) are.
Before going farther, here are some examples of differential k-forms, with
E = Rn .
For k = 0: a differential form of degree 0 is none other than a function
on E , of C p -class, with values in R.
For k = 1: let 1 be a differential 1-form. Then there exists an n-tuple
( f1 , . . . , fn ) of functions of C p -class such that
1 = f1 dx 1 + + fn dx n .
Differential forms
470
For k = 2: let 2 be a differential 2-form. There exist n(n 1)/2

functions ( f i j )1i< jn (of C p -class) such that
2 = f12 dx 1 dx 2 + f13 dx 1 dx 3 + + f1n dx 1 dx n
=
+ + fn1,n dx n1 dx n
1i< jn
f i j dx i dx j = f i j dx i dx j .
Note again the use of the convention of partial summation of repeated

indices.
Remark 17.29 Warning: in what follows, as in most books, the point x in space
where the extorior form is evaluated will be omitted from the notation, and we will
often write instead of ( x), when no ambiguity exists. Thus, for a 2-form ,
( 1 , 2 ) will be the real number obtained by applying the exterior 2-form ( x) to
the vectors 1 and 2 in E .
17.2.b Exterior derivative

DEFINITION 17.30 Let = i1 ,...,ik ( x) dx i1 dx ik be a differential k-
form. Since the functions x 7 i1 ,...,ik ( x) are differentiable, the differentials

d i1 ,...,ik ( x) are exterior 1-forms. The exterior derivative of the k-form is
the differential form of degree k + 1 given by
d = d i1 ,...,ik ( x) dx i1 dx ik .
Example 17.31 Let = x dx + y dy + z dz be a differential 1-form on R3 . Then we have
d x =
x
x
x
dx +
dy +
dz,
x
y
z
and similarly for d y and dz . It follows that

y
y
y
x
x
x
d =
dx +
dy +
dz dx +
dx +
dy +
dz dy
x
y
z
x
y
z

z
z
z
+
dx +
dy +
dz dz
x
y
z

y
y
x
z
x
z
dx dy +
dy dz +
dz dx.
=
x
y
y
z
z
x
THEOREM 17.32 Let be a differential k-form of class at least C 2 . Then dd
0, in other words we have identically d2 0.
Proof. We have by definition d = di1 ,...,ik ( x) dx i1 dx ik , with

di1 ,...,ik =
i1 ,...,ik
x
dx .
Integration of differential forms
471
It follows that
dd =
2 i1 ,...,ik
x x
dx dx dx i1 dx ik .
But since the form is of C 2 -class, Schwarzs theorem on the commutativity of second
partial derivatives implies that
2 i1 ,...,ik
x x
2 i1 ,...,ik
and
x x
dx dx = dx dx ,
which shows that the preceding expression is identically zero.
17.3
The general theory of integration of differential forms is outside the scope
of this book. Hence the results we present are only intended to give an
idea of the flavor of this subject. The interested reader can turn for more
details to a book of differential geometry such as Cartan [17], Nakahara [68],
or Arnold [9]. Before discussing the general case, we mention two special
situations: the integration of n-forms on Rn and of 1-forms on any Rn .
DEFINITION 17.33 (Integral of a differential n-form) Let be a sufficiently
regular domain in Rn (so is an n-dimensional volume) and let be a
differential n-form. Then there exists a function f such that we can write
( x) = f ( x) dx 1 dx n ,
and the integral of on is defined to be
Z
Z
def
=
f ( x) dx 1 dx n .
Be careful that this special case of integration of forms of degree n on a

space of the same dimension is indeed very special.
The next special case is that of differential 1-forms. Before trying to integrate a 1-form on a domain, we must look again at what is the nature
of . At a point x Rn , the value ( x) of is an exterior 1-form, that
is, simply a linear form which associates a real number to a vector. Where
may we naturally find vectors? In differential geometry, vectors are defined in
the same manner that velocity vectors are defined in physics: one chooses a
curve t 7 (t) going through x and one looks at what is the velocity of
the moving point (t) when it passes at the point x.
So, consider a path : [0, 1] Rn , t 7 (t). Then for any t [0, 1],
the velocity vector (t) = d /dt provides, after being fed to (t) (the
Differential forms
472
George Stokes (18191903), English physicist and mathematician,

discovered the concept of uniform convergence, but his most important contributions were in physics, in particular in fluid mechanics, where he investigated the differential equations introduced by
Navier that now bear their names. The formula linking the circulation of a vector field and the flux of its curl is also due to him.
exterior form given

byevaluating at the point (t)) a real number, namely,
g(t) = (t) (t) ; this real number may be integrated along the path,
and the resulting integral will give the definition of the integral of on the
path :
DEFINITION 17.34 (Integral of a differential 1-form on a path)
Let be a path and = f i dx i a differential 1-form. The integral of on

is defined as the integral
Z
Z
Z1
dx i
dt,
=
f i dx i =
f i (t)
dt
where the path is parameterized by
: [0, 1] Rn ,

t 7 (t) = x 1 (t), . . . , x n (t) .
Example 17.35 In the plane R2 , with = f dx + g dy, we obtain
f dx + g dy =
f (s)
dx
ds +
ds
g (s)
dy
ds.
ds
This recovers the notion already seen (in complex analysis) of contour integration.
In the general case, how do we integrate a differential k-form? Most importantly, on what domain should one perform such an integral? Take a 2-form,
for instance. At every point of this integration domain, we need two vectors.
This is easily given if this domain is a smooth surface, for the two vectors can
then be a basis of the tangent plane. A differential 2-form must be integrated along
a 2-dimensional surface.
As an example, we consider a 2-form 2 in the space R3 . We denote it
= f x dy dz + f y dz dx + f z dx dy.
(Note that we have changed conventions somewhat, using dz dx instead of
dx dz; this only amounts to changing the sign of f y .) Let S be an oriented
473
surface. At any point x in S , we can perform the orthogonal projection

of the basis vectors of R3 to the tangent plane of S at x (this is a linear
operation). Hence three vectors u x , u y , and u z are obtained, attached to each
point x S ; they are not linearly independent since they belong to the twodimensional tangent plane to the surface. Consider then f z dx dy( u x , u y ),
which is a real number. This may be integrated over the whole surface. Note
that dx dy( u x , u y ) = n e z , and therefore we obtain
Z
Z
2
=
f x dy dz + f y dz dx + f z dx dy
S
S
Z
Z
f x dx + f y dy + f z dz =
f n,
=
S
where n is the normal vector to the surface. The first integral involves the
differential form , whereas the integral on the second line is an integral of
functions.
This is easy to memorize: it suffices to replace dx dy by dz and, similarly
dy dz by dx and dz dx by dy.
To generalize this result, it is necessary to study the behavior of a differential form during a change of coordinates. This is done in sidebar 6. One shows
that a differential k-form may be integrated on a surface of dimension k,
that is, there appears a duality2 between
differential 1-forms and paths (curves);
differential 2-forms and surfaces;
differential 3-forms and volumes (of dimension 3) or three-dimensional
surfaces in Rn ;
etc.
The following result is then particularly interesting:
THEOREM 17.36 (Stokes) Let be a smooth (k + 1)-dimensional domain with k-
dimensional boundary (possibly empty), and let be a differential k-form on E .

Then we have
Z
Z
=
d.
This formula is associated with many names, including Newton, Leibniz, Gauss, Green,
Ostrogradski, Stokes, and Poincar, but it is in general called the Stokes formula.
Example 17.37 The boundary of a path : [0, 1] Rn is made of only the points b = (1)
and a = (0) . Let us consider the case n = 1 and let be a differential 0-form, that is, a
2
The following is meant by duality: the meeting of a k-form and a k-surface gives a
number, just as the meeting of a linear form and a vector gives a number.
Differential forms
474
function = f : R R. Then the Stokes formula means that we have

Zb
Z
Z1

df
df =
(t) (t) dt =
f (s) ds = f (b ) f (a),
a
0 dt
which is of course known to the youngest children.
17.4
Poincar
es theorem
DEFINITION 17.38 (Exact, closed differential forms) A differential form is
closed if d = 0.
A differential k-form is exact if there exists a differential (k 1)-form
A such that = dA.
THEOREM 17.39 Any exact differential form is closed.
Proof. Indeed, if is exact, there exists a form A such that = dA, and we then
have d = ddA = 0.
Note that an exact differential 1-form is none other than a differential (of
a function). Indeed, if is an exact differential 1-form, there exists a 0-form,
that is, a function f , such that = d f .
THEOREM 17.40 (Schwarz) Let be a differential 1-form. If we write =
f i (x) dx i , then the following are equivalent

a) is closed (i.e., d = 0) ;
b) we have
fj
fi
=
j
x
xi
for any i, j {1, 2, 3}.
Proof. Assume first that is closed. Then d = 0 = k f i dx k dx i ; since dx k

dx i = dx k dx i dx i dx k , this shows that k f i = i fk .
Conversely, k f i = i fk obviously implies that d = 0.
Knowing whether the converse to Theorem 17.39 holds, that is, whether
any closed form is exact, turns out to be of paramount importance. The
answer depends in a crucial way on the shape and structure of the domain of
definition of the differential form.
DEFINITION 17.41 (Contractible open set) A contractible open set in Rn is
an open subset which may be continuously deformed to a point. Hence, is

contractible if and only if there exists a continuous map C : [0, 1]
and a point p such that
Poincars theorem
475
Henri Poincar (18541912) is one the most universal geniuses

in the whole history of mathematics; he studied and often
revolutionized most areas of the mathematics of his time, and
also mechanics, astronomy, and philosophy. He created topology,
used this new tool to investigate differential equations, flirted
with chaos, studied asymptotic expansions, and proved that the
time series used in astronomy are of this type. He also gave a
major contribution to the three-body problem.
a) C (x, 0) = x for all x ;
b) C (x, 1) = p for all x .

Example 17.42 A star-shaped open set is contractible. Here, an open set is star-shaped if
and only if, for any x , the line segment

n
o
def
[0, x] = x ; [0, 1]
is contained is . In particular, Rn is star-shaped and hence contractible. Here we can choose

p = 0 and the map C may be defined by
x Rn
[0, 1]
C (x, ) = (1 ) x.
Example 17.43 The surface of a torus is not contractible.
Rn \ {0} is not contractible (although it is simply connected as soon as n 2).
THEOREM 17.44 (Poincare) Let Rn be a contractible open set in R. Then
any closed differential form on is exact.
In one important special case, a more general condition than contractibility is sufficient to ensure that a closed form is exact. Indeed, for differential
1-forms, we have
THEOREM 17.45 Let Rn be a simply connected and connected open set.
Then any closed differential 1-form on is exact.
Remark 17.46 Be careful that this theorem is only valid for differential 1-forms. In
general, it is not sufficient for a closed differential form to be defined on a simply
connected set in order for it to be exact. We will see below (see page 480) a simple
consequence of this for the field created by a magnetic monopole.
Notice that any contractible subset is simply connected (it suffices to look
at the image of a path under the contraction C to see that it also contracts
to a single point), but the converse is not true.
Differential forms
476
Counterexample 17.47 Consider the differential 1-form
(x y) dx + (x + y) dy
x 2 + y2

in R2 . It is easy to check that d = 0 on R2 \ (0, 0) , using Schwarzs theorem, for instance.
However, is not exact. One may indeed integrate to get

1
x
f (x, y) = log(x 2 + y 2 ) arctan
+ Cnt ,
2
y
=
but this expression

is singular at y = 0 and does not

admit any continuous continuation
on R2 \ (0, 0) . This is possible because R2 \ (0, 0) is not simply connected.
17.5
Relations with vector calculus:
gradient, divergence, curl
17.5.a
Differential forms in dimension 3
Consider now the case of differential forms in three-dimensional space R3 .

These are
the 0-forms, that is, functions;
the 1-forms, which are linear combinations of dx, dy, and dz;
the 2-forms which are linear combinatons of dx dy, dy dz, and
dz dx;
the forms of degre 3, multiples of the volume form dx dy dz.
Let be a differential 1-form, so that
= f i dx i = f x dx + f y dy + f z dz.
Let now be a smooth closed path in R3 . Then this path bounds an oriented
surface (oriented using the orientation of ), denoted S . So we have = S .
Since the exterior derivative of is given by

f y fx
fz f y
fx fz
d =
dx dy+
dydz+
dzdx,
x
y
y
z
z
x
Stokess formula gives
d =
S
Relations with vector calculus: gradient, divergence, curl
477
or, putting f = ( f x , f y , f z ) and using the fact that the components of d are
the components of the curl curl f ,
I
ZZ
2
curl f d =
f d,
(17.1)
S
or in other words still, the flux of the curl of f through a surface is equal to
the circulation of f along the boundary of this surface (Green-Ostrogradski
formula).
Similarly, if is a differential 2-form given by
def
= f x dy dz + f y dz dx + f z dx dy,
its exterior derivative is then

fx f y fz
d =
+
+
dx dy dz = (div f ) dx dy dz,
x
y
z
and if S is a closed surface bounding a volume V such that S = V , we

obtain
Z
Z
ZZ
ZZZ
2
=
d,
that is,
f d =
div f d3 r.
V
Finally, if = f is a differential 0-form, that is, a function, since the

differential d f satisfies d f ( n) = grad f n, the Stokes formula expresses that
Z

grad f d = f (1) f (0) .
Note that the relations just stated between exterior derivatives and the
gradient, divergence and curl operators show that the identity d2 0 may
simply be expressed, in the language of the man in the street, as
div curl 0
and
curl grad 0.
17.5.b Existence of the scalar electrostatic potential

Let f = ( f1 , f2 , f3 ) be a vector field on R3 , and assume it satisfies the condition
f = curl f 0
(17.2)
3
at every point of a domain R which is simply connected. Define a
differential 1-form on by
= f i dx i = f x dx + f y dy + f z dz.
As we have seen, the exterior derivative of is

f y fx
fz f y
fx fz
d =
dx dy+
dydz+
dzdx,
x
y
y
z
z
x
478
Differential forms
hence condition (17.2) can be expressed more simply by stating that

d 0.
(17.2 )
Poincars theorem (the second version, Theorem 17.45) ensures that there
exists a differential 0-form, that is, a function : R R, such that
= d =
dx +
dy +
dz.
x
y
z
Comparing with the definition = f x dx + f y dy + f z dz, it follows that

f = .
This proves the following theorem:
THEOREM 17.48 (Existence of the scalar potential) Let R3 be a simply
connected open set in R3 , and let f : R3 be an irrotational vector field on ,

that is, such that curl f 0. Then there exists a potential : R such that
f = .
Application to electrostatic
Since the electrostatic field satisfies E 0 on R3 , both with and without
charges3 it follows that it derives from a potential:
E = .
It is important to remark that the scalar potential only exists because
R3 is simply connected. If the equation were only satisfied in a non-simply
connected subset R3 , there would be no reason for the potential to exist
globally (its local existence being of course always true).
Too many physics books do not hesitate to state things in roughly the
following way: since we know that curl grad 0, it follows that: if curl E 0,
then E = grad . This, we just saw, is a mathematical absurdity and, worse,
a dreadful logical mistake.4 Moreover this may discredit physics in the eyes of
students who, in general, are perfectly aware that they are being hoodwinked,
even if they may not know the necessary tools and terminology to explain
the problem.5 It is perfectly possible to remain rigorous in physics; students
should be made aware of this.
Should you, reader, teach physics one day, do not make this mistake!
3
This is one of the homogeneous Maxwell equations, giving the global structure of fields in
the setting of electrostatics. Hence all magnetic phenomena have been neglected.
4
This amounts to mixing up the logical statements (P Q ) and ( Q P )!
5
And who, moreover, would be quite incapable of finding a counterexample: it is true that
E comes from a potential, but for a very different reason.
Relations with vector calculus: gradient, divergence, curl
479
17.5.c Existence of the vector potential

Let now f : R3 be a vector field with vanishing divergence:
div f = f 0
on .
(17.3)
Define a differential 2-form

def
= f1 dx 2 dx 3 + f2 dx 3 dx 1 + f3 dx 1 dx 2 .
Then the exterior derivative of is simply

f1
f2
f3
d =
+
+
dx 1 dx 2 dx 3 ,
x1 x2 x3
and the condition (17.3) can be expressed again simply as
d = 0.
Poincars theorem (Theorem 17.44 now) ensures that if is contractible, there
exists a differential 1-form A = A1 dx 1 + A2 dx 2 + A3 dx 3 such that
= dA

A2 A1
A3 A2
1
2
=
dx dx +
dx 2 dx 3
x2
x2
x3
x1

A1 A3
+
dx 3 dx 1 ,
x3
x1
which, putting A = (A1 , A2 , A3 ) and comparing with the definition of ,
shows that
f = A = curl A.
THEOREM 17.49 (Existence of the vector potential) Let R3 be a contrac-
tible open subset and let f : R3 be a vector field with zero divergence on :
div f 0. Then there exists a vector potential A : R3 such that f = A.
In the space R3 , the magnetic field satisfies the relation div B 0, and
this independently of the distribution of charges. It follows that the magnetic
field, on R3 which is contractible, comes from a vector potential A.
Remark 17.50 In the general setting of electromagnetism, one first proves the existence of the
vector potential A; then, from the equation curl E = B/ t, it follows that curl(E +
A/ t) = 0 and hence the existence of the scalar potential such that
E =
on any contractible open set.
A
t
Differential forms
480
17.5.d Magnetic monopoles

The magnetic field B satisfies div B 0, and so it comes from a vector potential; or, more precisely, this is ensured by Poincars theorem when considering
a contractible set (in particular it is not sufficient to have a simply connected
set).
What happens if we consider a magnetic monopole, that is, a point-like particle
which transforms one the Maxwell equations to be
div B( r , t) = qm ( r ) ?
(17.4)
The relation div B 0 is satisfied now only on R3 \{0}, which is not contractible.
It is possible to show then that there does not exist a vector potential A defined
on R3 \ {0} such that we have B = A at every point of R3 \ {0}. But
since a local description of the magnetic field using a vector potential remains
possible, it is possible to show that one may defined two potentials A1 and
A2 defined on a neighborhood of the upper half-space (resp. lower half-space)
such that the relations A i = B for i = 1, 2 hold on the respective domains
where the vectors are defined. For instance, A1 and A2 may be defined on
the following grayed domains:
0
A1
0
A2
Hence the existence of magnetic monopoles would have as a consequence

that the vector potential does not exist. Experiences have been attempted
to discover magnetic monopoles (which have amusing theoretical properties;
for instance, their existence automatically implies the quantification of the
elementary electric charge), but without success.
17.6
Electromagnetism in the language
of differential forms
Physicists see the vector potential usually as a vector, but it should rather
be seen as a differential form. Indeed, if A itself does not have an intrinsic
meaning,H the integral of A along a path has one: for a closed path , the
integral A d is equal to the flux of the magnetic field through a surface
S bounded by (see formula (17.1)).
Electromagnetism in the language of differential forms
481
This is why, in special relativity, the components of the potential fourvector (vector potential + scalar potential) are denoted in contravariant notation:
A = A dx .
The Faraday tensor (which is a differential 2-form or equivalently a field
0
of 2 -tensors) is defined as the exterior derivative of A:
F = dA, namely, F = 12 F dx dx = 12 A A dx dx ,
where complete summation on repeated indices is performed. Denoting the

potential four-vector A = (A0 , A), E = grad A0 A/ t and B = curl A,
the coordinates of F are given by
F01 = E x ,
F02 = E y ,
F12 = B z ,
F13 = B y ,
F03 = E z ,
F23 = B x ,
that is, we have

F = E x dx dt + E y dy dt + E z dz dt B z dx dy B x dy dz + B y dz dx.
Since F = dA, its exterior derivative dF is zero, which means that

B x B y Bz
0=
+
+
dx dy dz
x
y
z

Ez E y Bx
+
+
dt dy dz
y
z
t

E x Ez B y
+
dt dz dx
+
z
x
t

E y E x Bz
+
dt dx dy.
+
x
y
t
If we reorder the terms, we obtain the Maxwell equations without source:
curl E = B Maxwell equations

dF = 0
t
without source.
div B = 0
How is it possible to obtain the Maxwell equations with source? We need

to involve derivatives of E and B, but we cannot get this from dF (which is
identically zero). Another differential form is needed, constructed using E
and B and involving derivatives of the first order.
Now notice that the space of exterior 1-forms is a vector space of dimen4
sion 1 = 4 (generated by dt, dx, dy, and dz), and the the space of exterior
Differential forms
482
4
3-forms is of dimension 3 = 4 also. Using the Hodge operator, an isomorphism between both spaces is constructed:
DEFINITION 17.51 (Hodge operator) The Hodge operator, denoted , is a

linear map defined on the spaces of exterior forms on Rn :
k (Rn ) (nk) (Rn ),

1
dx i1 dx ik 7
dx ik+1 dx in
(n k)! i1 ,...,ik ,ik+1 ,...,in
with complete summation on repeated indices. For any k [[0, n]], the Hodge
operator is an isomorphism between the space of k-forms and the space of
(n k)-forms.
Example 17.52 Consider the case n = 3. Then dx =
Similarly
(dx dy) =
1
(31)!

dy dz dz dy = dy dz.
1
dz = dz.
(3 2)!
The Hodge dual of a k-form is thus a (3 k)-form.

In the case n = 4, we find for instance that dx = dy dz dt and that (dx dy) = dz dt.
A k-form is mapped to a (4 k)-form.
From the differential 2-form F , we can thus construct another 2-form,

namely, F . It is given by
F
= B x dx dtB y dydtB z dzdt+E z dx dy+E x dydz+E y dzdx.
Let J be the differential 1-form J = dt + j x dx + j y dy + jz dz, and consider

its dual (image) by the Hodge operator
J = dx dy dz j x dt dy dz j y dt dz dx jz dt dx dy.
The equation
d F = J
(17.5)
is then equivalent to
E
=j
t
div E =
curl B
Maxwell equations
with source.
Before concluding, a last remark about the differential form J . The

equation (17.5) means that this form is exact (being the exterior derivative
of F ). Hence its exterior derivative is identically zero, that is, we have

j x j y jz
d J =
+
+
+
dt dx dy dz 0,
t
x
y
z
Electromagnetism in the language of differential forms
483
which is simply the equation of conservation of charge, which is indeed known

to follow from the Maxwell equations.
Note that the components of F and F may be written down in matrix
form. They are
0
Ex
Ey
Ez
E
0 B z
By
,
F = (F ) = x
Bz
0 B x
E y
E z B y
Bx
0
0 B x B y B z
0
E z E y
F = ( F ) = B x
.
B E
0
Ex
y
z
Bz
E y E x
0
The coordinates of F and F satisfy also

F
1
= F ,
2
where is the totally antisymmetric tensor defined by 1234 = 1, which is

called the Levi-Civita tensor. It satisfies = 1 if (, , , ) is an
even permutation of (1, 2, 3, 4), = 1 if it is an odd permutation,
and = 0 if any index is repeated more than once.
Remark 17.53 What are the laws of electromagnetism in two dimensions? If we want to keep
the local nature of the fundamental laws of the universe, it is their differential expression which
should be preserved in order to extend them to such imaginary universes. Thus, one should
keep the structure of the Faraday tensor for the electromagnetic field. This may be written here
0
Ex
Ey
0
B ,
F = E x
E y
B
0
which shows that the electric field has two components, whereas the magnetic field is a scalar
field (or rather a pseudo-scalar, which changes sign during a change of basis that reverses the
orientation).
484
Differential forms
PROBLEM
Problem 6 (Proca lagrangian and Cavendish experiment) We want to change the Maxwell
equations in order to describe a photon with nonzero mass. For this purpose, we start
from the lagrangian description of the electromagnetic field, and add a mass term to the
lagrangian. Then we look for detectable effects of the resulting equations.
I. Recall that a lagrangian density describing the laws6 of electromagnetism is given by

1
def
L = F F J A /0 ,
4
with units where h} = c = 1 [49, 50, 57], where the F are the covariant coordinates
of the Faraday tensor, related to the potential four-vector by F = A A . We
have denoted here = / x and used the Einstein convention for the summation
of repeated indices (see Chapter 16 concerning tensors).
1. Show that the equations of movement, given by the Euler-Lagrange equations
L
L
=
A
( A )
provide the Maxwell equations with source. Write them in Lorenz gauge, in which
the condition A = A = 0 is imposed.
2. A mass term of the type
2 2 2
A =
A A ,
2
2
is added to the classical lagrangian to obtain the Proca lagrangian
2
1
def
A A .
LP = F F J A /0 +
4
2
Write the equations of movement associated with this lagrangian. Show that, in
Lorenz gauge, they are given by
( + 2 )A = J /0 .
3. Express the relation between the frequency of an electromagnetic wave and its wave
number. Show that has the dimension of mass (still with h} = c = 1). Using the
de Broglie relations, show that the waves thus described are massive.
4. Consider the field created by a single, immobile, point-like particle situated at the
origin. Show that only the component A0 = is nonzero. Show that this scalar
potential is of the Yukawa potential type and is equal to
q
e r
( r) =
.
40
r
What is the typical decay distance for this potential?
II. In 1772, Cavendish realized an experiment to check the 1/r decay of the Coulomb
potential [50].
Consider two hollow spheres, perfectly conducting, concentric with respective radii R1
and R2 , with R1 < R2 . At the beginning of the experiment, both spheres are electrically
neutral, and the exterior sphere is given a potential V . Then the spheres are joined
(with a thin metallic wire for instance).
6
That is, the Maxwell equations for the electromagnetic field, and the Newton equation
with Lorentz force for particles.
485
1. Show that, if the electrostatic potential created by a point-like charge follows the
Coulomb law in 1/r , the inner sphere remains electrically neutral.
2. We assume now that the electromagnetic field is described by the Proca lagrangian
and that the electric potential created by a point-like charge follows the Yukawa
potential. Show that this potential satisfies the differential equation
( +2 ) = /0 .
3. Show that the electric field has a discontinuity /0 when passing through a surface
density of charge .
4. Compute the potential at any point of space after the experiment is ended (both
spheres are at the same potential V ). Deduce what is the charge of the inner sphere.
(One can assume that is small compared to 1/R1 and 1/R2 and simplify the
expressions accordingly.)
5. Knowing that Cavendish, in his time, could not measure any charge on the inner
sphere, and that he was able to measure charges as small as 109 C, deduce an upper
bound for the mass of a photon. Take V = 10 000 V and R1 R2 = 30 cm.
Any comments?
SOLUTION
I. 1. In Lorenz gauge, one finds F = J , which is indeed the covariant form of the
Maxwell equations.
2. Immediate computation.
3. Notice first that, in analogy with the Klein-Gordon equation, the additional term
deserves to be called a mass term. If we consider a monochromatic wave with
frequency and wave vector k, the equation of movement in vacuum (+2 ) A =
0 gives 2 + k 2 + 2 = 0. The de Broglie relations link frequency and energy on
the one hand, and momentum and wave vector on the other hand. We obtain then
E 2 = p 2 + 2 , which may be interpreted as the Einstein relation between energy,
momentum, and mass, taking as the mass. In fact, it is rather h}/c which has
the dimension of a mass.
4. We have J = (q, 0); the components Ai for i = 1, 2, 3 therefore are solutions of
the equation ( + 2 ) Ji 0 and are zero up to a free term.
The component = A0 , on the other hand, satisfies
d
+ 2 = q ( r )/0 ,
dt
which is independent of time. Looking for a stationary solution, we write therefore
( +2 ) = q ( r)/0 ,
which (see Problem 5, page 325) can be solved using the Fourier transform and a
residue computation, and yields the formula stated. The typical decay distance is
therefore 1 .
II. 1. When the exterior sphere is charged with the potential V , the potential at any point
exterior to the two spheres is (using for instance Gausss theorem on the electric
field)
R
( r) V 2
for all r > R2 .
r
Differential forms
486
The inner potential, on the other hand, is uniform and equal to V (same method).
In particular, the inner sphere is already at the same potential V . The electric wire
is also itself entirely at the potential V . Hence there is no electric field, no current
is formed, and the inner sphere remains neutral.
2. If electromagnetism follows the Proca lagrangian, then the Poisson law is not valid
anymore, and consequently, Gausss theorem is not either. The potential created by
a point-like charge being
( r ) =
1
e r ,
40 r
it satisfies the differential equation (2 ) = /0 (which is indeed the differential equation of Question I.2 in stationary regime); by convolution, the potential
created by a distribution of charge therefore satisfies (2 ) = /0 (recall
that ( ) = () ).
3. Gausss theorem is not valid anymore, since we have
div E = = /0 2 .
Consider a charged surface element, surrounded by a surface S .
S
Then we have
ZZ
E dS =
ZZZ
d3 r =
ZZZ
2
0
d3 r .
But, close to , the Yukawa

potential is regular (i.e., behaves like the Coulomb
RRR
potential), and hence V 2 d3 r tends to 0 when the volume tends to 0. The
distribution is a Dirac surface distribution surfacique. It follows that the discontinuity of the normal component of the electric field is /0 .
4. Since the potential has radial symmetry, it satisfies the equation
1 2
(r ) 2 = 0
r r2
at every point in space except where r = R1 or r = R2 . The general form of the
potential is therefore
A
B
(r ) = cosh r + sinh r.
r
r
Denote by , , and the functions giving the potential (respectively) inside the
inner sphere, between the two spheres, and outside the exterior sphere.
The potential is bounded in the neighborhood of the origin, and therefore may
be written
sinh(r )
(r ) = A1
,
r
which, together with the condition of continuity of the potential (R1 ) = V , gives
(r ) = V
R1 sinh(r )
.
r sinh(R1 )
487
The potential is of the form

B
A
cosh r + sinh r
for all R1 < r < R2 ,
r
r
and moreover satisfies (R1 ) = (R2 ) = V , which leads to the system of equations
A e R1 + B e R1 = R1 V ,
A e R2 + B e R2 = R2 V .
(r ) =
Solving this linear system, we obtain

(r ) =

V
R sinh (R2 r ) + R2 sinh (r R1 )
r sinh(R) 1
with R = R2 R1 . Then, differentiating, we get the values of the electric field on

the left and on the right of R1 , and the difference is

E (R1+ ) E (R1 ) =
r r=R1
r r=R1

R2
= V coth R1 + coth R
= .
R1 sinh R
0
Thus, for small , we have
=
2 V 0 R2
(R + R2 ).
6 R1 1
The total charge is therefore

Q = 4R12 =
2 2
V R1 R2 0 (R1 + R2 ).
3
5. With the data in the text, since Q 109 C, we obtain 1 6 m, or, adding
correctly the proper factors h} and c ,
3. 1043 kg,
which is a very remarkable result. (Compare, for instance, with the electron mass
me = 9.1 1031 kg.)
Differential forms
488
Sidebar 6 (Integration of differential forms) Let be a differential k-form.

In Rk , it would be easy to give a meaning to the integral of on (a subset of) Rk
(see Definition 17.33, page 471).
In order to define the integral of on a k-dimensional surface S (also called
a k-dimensional manifold), one starts by transforming S into a manifold ,
which is indeed a subset of Rk , by means of a change of coordinates. Let f denote
the inverse change of coordinates, that is, the map from Rk into S :
f : S Rn ,
x 7 f ( x).
A point x Rk in the space Rk is mapped to f ( x). How are the tangent

vectors in Rk transformed? Since tangent vectors correspond to differences between
two points of Rk , the image by the change of coordinates of a vector v attached
to the point x will be v 7 d f x . v, that is, the result of evaluating the differential
form d f x (given by the differential of f at the point x) at the vector v. It is often
customary to write f = d f the differential of f (without indicating precisely at
which point it is considered).
From the change of coordinates, a k-differential form on , the inverse image
of by f , which is denoted f , may be defined by putting
( f ) :
(Rk )k R,
( v1 , . . . , vk ) 7 f ( v 1 , . . . , vk ) = ( f v1 , . . . , f vk ).
def
Using the definition of the integral of a differential k-form on Rk , it is then possible

to define:
Z
Z
f .
def
Of course, this definition can be shown to be independent of the choice of f .

f
f
S
Rk
Rn
Chapter
18
Groups
and group representations
18.1
Groups
DEFINITION 18.1 A group (G , ) is a set G with a product law defined on
G G , such that
1. is associative: (g h) k = g (h k) for all g, h, k G ;
2. (G , ) has an identity element e, that is, an element of G such that

e g = g e = g for any g G ;
3. any element in G is invertible, that is, for any g G , there exists an
element h G such that g h = h g = e. This element is unique; it is
denoted g 1 and is called the inverse of g.
DEFINITION 18.2 A group (G , ) is abelian or commutative if the product
law is commutative, that is, if g h = h g for any g, h G . An abelian group

is often denoted additively: one writes g + h instead of g h.
Groups and group representations
490
For commutative groups, the definition can thus be rephrased as follows:

DEFINITION 18.3 An additive group is a pair (G , +) such that
1 . + is an addition law on G , which is associative and commutative;

2 . (G , +) has an identity element 0G such that
g G
0G + g = g + 0G = g;
3 . any G is invertible, that is, for any g G , there exists an element

h G such that g + h = h + g = 0G . This element is denoted g and
is called the opposite of g.
Example 18.4 The set GLn (R) of invertible square matrices with real coefficients of size n, with
the law given by the product of matrices, is a group; for n 2, it is non-abelian.
Example 18.5 The set of symmetries of a molecule, with the law given by composition of
maps, is a group.
Example 18.6 The set of rotations of the vector space R2 (rotations fixing the origin) is a
commutative group.
Example 18.7 Consider the oriented space E = R3 with its euclidean structure. Let G denote
the set of rotations of E centered at the origin 0. Then (G , ) is a noncommutative group.

Indeed, if R x denotes the rotation with axis O x and angle /2, and if R y denotes the rotation
with axis Oy and angle /2, the following transformations of a domino show that R x R y 6=
R y Rx :
z
y
Rx
Ry
Ry
Rx
z
y
x
Linear representations of groups
491
18.2
Linear representations of groups
A group may be seen as an abstract set of objects, together with a table
giving the result of the product for every pair of elements (a multiplication
table). From this point of view, two groups (G , ) and (G , ) may well have
the same abstract multiplication table: this means there exists a bijective map
: G G which preserves the product laws, that is, such that for any g1 , g2
in G , the image by of the product g1 g2 G is the product of the respective
images of each argument:
(g1 g2 ) = (g1 ) (g2 ).
In such a case, the groups are called isomorphic.
DEFINITION 18.8 A map that preserves in this manner the group structure
is called an homomorphism or a group morphism; if it is bijective, it is a

group isomorphism.
Example 18.9 The set M of matrices
Mx =
1
0
x
,
1
with the product given by the product of matrices, is a group. Since M x M y = M x+ y for any
x, y R, it follows that the (obviously bijective) map : x 7 M x is a group isomorphism
between (R, +) and (M , ).
Alternately, one speaks of a representation of the group G by G if there

exists a homomorphism : G G (not necessarily bijective).
Example 18.10 Let U be the set of complex numbers z such that |z| = 1. The map : R U
defined by () = e i is a group morphism between (R, +) and (U, ), but it is not an

isomorphism. Its kernel, defined as the set of elements with image equal to the identity
element in U, is given by
Ker = { R ; () = 1} = 2Z.
The map : R U is a (nonbijective) continuous representation1 of the additive group R.
(Note that the groups R and U have very different topological properties; for instance, R is
simply connected, whereas U is not.)
On the other hand, if we denote by G the group of rotations of R2 centered at the
origin 0, the map : G U which associates e i to a rotation with angle is indeed a group
isomorphism. Moreover, it is continuous. It follows that G , like U , is not simply connected.
By far the most useful representations are those which map a group (G , )
to the group of automorphisms of a vector space or, equivalently, to the group
GLn (K) of invertible square matrices of size n:
1
To speak of continuity, we must have a topology on each group, which here is simply the
usual topology of R and of U C.
492
DEFINITION 18.11 Let K be either R or C. A linear representation of a
group G defined over K is a pair (V , ), where V is a K-vector space and

is a homomorphism of G into G(V ), the group of linear automorphisms
of V.
Such a representation is of finite degree if V is a finite-dimensional vector
space, and the degree of the representation is the dimension n of V . In this
case, G(V ) is isomorphic to GLn (K), and we have a matrix representation.
By abuse of notation, one sometimes refers to V only, or to only, as a
representation of the group G .
DEFINITION 18.12 Let (V , ) be a linear representation of a G . This repre-
sentation is faithful if is injective, that is, if distinct elements of the group

have distinct images (this means no information is lost by going through the
map ).
On the other hand, the representation is trivial if (g) = Id V for any
g G , i.e., if we have (g)(v) = v for all g G and all v V (which means
all information in G is lost after applying ).
Example 18.13 Consider Example 18.9 again. The map : R GL2 (R) which associates the
matrice M x to a real number x is a faithful matrix representation of R.
18.3
Vectors and the group S O(3)
In this section, physical space is identified with the vector space E = R3 .
DEFINITION 18.14 The special orthogonal group or group of rotations of
euclidean space E = R3 , denoted S O(E), is the group of linear maps R :

R3 R3 such that
the scalar product is invariant under R:
a, b R3

( a| b) = R( a)R( b) ;
R preserves orientation in space: the image of any basis which is positively oriented is also positively oriented.
The rotation group SO(E ) will be identified without further comment with
the group S O(3) of matrices representing rotations in a fixed orthonormal
basis.
Denoting by R the matrix representing a rotation R in the (orthonormal)
canonical basis of E , we have
( a| b) = (R a|R b) = ( a|t RR b)
for all a, b E
Vectors and the group SO(3)
493
(using the invariant of the scalar product under R), which shows that R is an
invertible matrix and that
t R = R 1 .
(18.1)
A matrix for which (18.1) holds is called an orthogonal matrix.
Thus linear maps preserving the scalar product correspond to orthogonal
matrices, and they form a group, called the orthogonal group, denoted O(E ).
The group of orthogonal matrices is itself denoted O(3). Note that the matrix
I3 , although it is orthogonal, is not a rotation matrix; indeed, it reverses
orientation, and its determinant is equal to 1.
The group O(3) is therefore larger than SO(3). More precisely, we have:
THEOREM 18.15 The group O(3) is not connected; it is the union of two connected
subsets, namely, SO(3) and {R ; R SO(3)}, on which the determinant is equal,

respectively, to 1 and 1.
Since the identity matrix lies in SO(3), it is natural to state that SO(3) is
the connected component of the identity in O(3).
To summarize, the notation SO(3) means
S for special (i.e., with determinant 1),

O for orthogonal (such that t R = R 1 ).
Consider then a rotation R , characterized by a vector R3 with norm

[0, 2[ (the angle of the rotation), directed along a unit vector n (the axis
of the rotation): = n.
Denote by (e1 , e2 , e3 ) the basis vectors in space and by (e1 , e2 , e3 ) their
images by a rotation R . With matrix notation (it is easy to come back to
tensor notation) we denote by R(
) the matrix defined by the relations2
ei =
3
X
j=1

R(
) i j e j .
2
Note : the matrix R(
) is not the matrix for the change of basis L defined Equation (16.1)
on page 434, but its transpose. Here, we have, symbolically,

e
e
1
1
) e2 ,
e2 = R(
e3
e3
which means that one can read the coordinates of the new vectors in the old basis in the rows
of the matrix (rather than in the columns). The matrix linking the coordinates of a (fixed)
point in both bases is R(
)1 = t R(
) :

x
x
t

) y .
y = R(
z
z
494
A rotation R( e z ) with angle around the axis O z is represented by the

matrix
cos sin 0
def
R z () = R( ez ) = sin cos 0
0
0
1
Similarly, the matrices of rotations around the axis O x and O y are defined
by
1
0
0
def
sin
R x ( ) = R( e x ) = 0 cos
0 sin
cos
and (be careful with the sign in this last formula!)
cos 0 sin
def
1
0 .
R y () = R( e y ) = 0
sin 0 cos
Remark 18.16 When speaking of space rotations, there are two different points of view, for
a physicist. The first point of view (called active) is one where the axes of the frame (the
observer) are immobile, whereas the physical system is rotated. In the second point of view,
called passive, the same system is observed by two different observers, the axes of whose
respective frames differ from each other by a rotation. In all the rest of the text, we only consider
passive transformations.
For a rotation R(
) parameterized by = n, we can express the vector
in terms of the canonical basis:
= e1 + e2 + e3 .
How does one deduce the representation for the rotation R(
)? It is tempting
to try the product R x ( ) R y () R z (); however, performing the product in
this order seems an arbitrary decision, and the example of the domino on
page 490 should convince the reader that this is unlikely to be the correct
solution. Indeed, in neither the first nor the second line is the
p final state of
the domino the same as its state after a rotation of angle / 2 around the
first diagonal in the O x y-plane.
A proper solution is in fact quite involved. It is first required to use
infinitesimal transformations, that is, to consider rotations with a very small
angle . Then it is necessary to explain how non-infinitesimal rotations may
be reconstructed from infinitesimal transformations, by means of a process
called exponentiation (which is of course related to the usual exponential
function on R).
Consider then a rotation with axis O z and infinitesimal angle : it is
given by
1 0
R( ez ) = 1 0 = I3 + Jz
0
0 1
Vectors and the group SO(3)

where we put
Similarly, define
and
495
0 1 0
= 1 0 0 ,
=0
0 0 0

def dR( e z )

Jz =
d

0 0

dR(
e
)
def
x
Jx =
=
0 0
d =0
0 1
1
0
0 0 1
= 0 0 0 .
=0
1 0 0

def dR( e y )

Jy =
d
DEFINITION 18.17 Denote J1 = J x , J2 = J y , and J3 = J z . The matrices J i
thus defined are the infinitesimal generators of rotations. We denote by J

the vector J = ( J1 , J2 , J3 ).
How does one now pass from infinitesimal transformations to finite
transformations?3 A first method is to remark that the set of matrices of the
form R( e z ) is a representation of the additive group
(R, +), what is called a
one-parameter group,4 since we have R ( + ) e z = R( e z ) R( e z ). An
explicit computation shows that

sin 0
cos 0 0
dR( ez )
= cos 0 sin 0 0 = Jz R(0 e z ),

d
=0
0
0
1
and hence the matrix-valued function R( ez ) is a solution of the ordinary

differential equation with constant coefficients
M = Jz M
for which the solutions are given by
2 2
n n
Jz + +
J + .
2!
n! z
(This series is absolutely convergent, and therefore convergent.) This is a
general fact, and the following theorem holds:

THEOREM 18.18 Let R() R be a one-parameter group of matrices. Then its
infinitesimal generator given by J = R (0) satisfies R() = e J for any R.
R( ez ) = R(0)exp( Jz ) = exp( Jz ) = I2 + Jz +
Finite in the usual physicists sense, meaning not infinitesimal.

I.e., a group with elements of the type g x , x running over the group of real numbers R,
such that g y g x = g y+x and gx = g x1 for any x, y R.
4
496

Proof. Indeed, we have R(0) = I and
R( + h) R()
dR()
= lim
= lim
h0
h0
d
h
R(h) Id
h
R() = J R(),
which shows that R() is the solution of the differential equation with constant coefficient R = J R such that R(0) = I , that is, R() = e J .
THEOREM 18.19 Let R be a rotation characterized by the vector = n (angle ,
axis directed by the vector n). Then, denoting as before J = ( J1 , J2 , J3 ), the matrix
) is given by
representing R(
R(
) = exp(
J) = exp(1 J1 + 2 J2 + 3 J3 ).
It is important here to be very careful that the matrices Ji do not commute

with each other, so that this expression is not a rotation with angle 1 around
the x-axis, followed by a rotation with angle 2 around the y-axis, and a
rotation with angle 3 around the z-axis:
exp(1 J1 + 2 J2 + 3 J3 ) 6 exp(1 J1 ) exp(2 J2 ) exp(3 J3 ).
However, using the exponentials, the matrices Ji can provide any rotation
matrix. Even better: it turns out that, in practice, it is only necessary to know
the commutation relations between the matrices Ji . The reader will easily check
that they are given by
[ J1 , J2 ] = J3 ,
[ J2 , J3 ] = J1 ,
and
[ J3 , J1 ] = J2
(recall that [A , B] = A B B A for any two square matrices A and B), which
can be summarized neatly by the single formula
[ Ji , J j ] = i jk Jk ,
(18.2)
where summation over the repeated index (i.e., over k) is implicit, and where
the tensor is the totally anti-symmetric Levi-Civita tensor. The commutation
relations (18.2) define what is called the Lie algebra of S O(3). The coefficients
i jk are the structure constants of the Lie algebra of SO(3).
Remark 18.20 The minus sign in the formula is not important; if we had been dealing with
active transformations instead of passive ones, the matrices Jk would have carried a minus sign,
and the commutation relations also.
Knowing the Lie algebra of the group (i.e., knowing the structure constants)
is enough to recover the local structure of the group. (However, we will see
that there exist groups which are globally different but have the same Lie
algebra, for instance, SO(3) and SU(2), or more simply SO(3) and O(3); but
this is not a coincidence, since SU(2) is a covering of SO(3).)
Remark 18.21 Since SO(3) is a group, performing a rotation with parameter followed by a
rotation with parameter is another rotation, with parameter . However, there is no simple
relation between , and .
The group SU(2) and spinors
497
DEFINITION 18.22 Objects that are transformed by a matrix of SO(3) during
a passive rotation of coordinate axes are called vectors. It is sometimes said

that SO(3) is the vector representation of the group of rotations.
Example 18.23 Consider a rotation R of R3 . The differential of R at any point is dR = R,
which shows that velocity vectors transform according to R. In other words velocity vectors
are indeed objects that transform according to a vector representation.
A vector, in the representation-theoretic sense here, is therefore more than

simply an element of a vector space: it is an object generating the representation SO(3) of the rotation group [18, 28].
18.4
The group S U(2) and spinors
Consider now the complex vector space E = C2 , with its canonical
Hilbert structure, that is, with the hermitian scalar product

x
y
( x| y) = t x y = x1 y1 + x2 y2 ,
where x = 1 and y = 1 .
x2
y2
DEFINITION 18.24 The special unitary group of E , denoted S U(E ), is the
group of linear maps from E to itself which preserve the scalar product and
have determinant one.
If M is the matrix representing such a linear map in the canonical basis,

we have therefore:

( x| y) = (M x|M y) = t (M x) M y = x t M M y
for any x, y E , which implies that t M M = I2 .
DEFINITION 18.25 The conjugate transpose of a matrix M is denoted

def
M = t M.
A matrix such that M M = I2 only also satisfies
|det M |2 = det M det t M = det(M M ) = 1.
Such a matrix is a unitary matrix. Its determinant is of the form e i for

some real number .
Example 18.26 The matrix M = diag(e i , e i ) is unitary, but does not belong to SU(2) if
+ 6 0 mod 2.
498
DEFINITION 18.27 The group of complex 2 by 2 matrices such that
is denoted S U(2).
M M = I2
and
det M = 1
The notation SU(2) comes from
S for special (i.e., with determinant 1),

U for unitary (i.e., with t M = M 1 ).
Recall that any rotation matrix could be expressed as the exponential of a

linear combination of elementary matrices Jk ; a similar phenomenon holds
for matrices in SU(2). First, notice that any matrix in SU(2) is of the form

U=
with ||2 + ||2 = 1.

Exercise 18.1 Prove this property.
If we write = a + ib and = c + id , it follows therefore that a 2 +

+ c 2 + d 2 = 1. We can write a = cos(/2) for some and then write
(b, c , d ) = sin(/2) n, where the real vector n = (n x , n y , n z ) is of norm one.
Then U is given by

U (
) = cos
I2 + i sin
n ,
2
2
b2
where
def
= (1 , 2 , 3 ),
the matrices i being none other than the Pauli matrices, well known to
physicists:
0 1
0 i
1 0
1 =
2 =
3 =
.
1 0
i 0
0 1
Those matrices have trace zero (they are traceless in the physicists language)
and are hermitian, which means that we have
i = i
for i = 1, 2, 3.
The Pauli matrices play for SU(2) the same role that the Jk play for SO(3), as
we will see. First, notice that the reasoning above can be reversed:
PROPOSITION 18.28 Let = n be a vector in R3 . The matrix
is hermitian,
and the matrix U (

) = e i /2 is unitary. Moreover, we have
U (
) = (cos /2) I2 + (i sin /2) n
=
cos 2 + in z sin 2
(in x + n y ) sin 2
(in x n y ) sin 2
cos 2 in z sin 2
499
Proof. A direct computation shows that U (

) t U (
) = I2 .
PROPOSITION 18.29 (Properties of the Pauli matrices)
Commutation relations:
Anticommutation relations:
i i i j
,
2 2
= i jk
ik
;
2
i j + j i = 2 i j I2 ,
i j = i j I2 + i i jk k ,
tr ( i ) = 0,
i = i .
The groups SO(3) and SU(2) are quite similar. Indeed, for any vector
R3 , we construct two matrices
R(
) = exp(
J) SO(3)
and
U (
) = exp(i
/2) SU(2),
using the matrices ( J1 , J2 , J3 ) and (1 , 2 , 3 ), respectively, which satisfy the

same commutation relations:

i i i j
i
[ Ji , J j ] = i jk Jk
and
,
= i jk k .
2 2
2
The structure constants of SO(3) and SU(2) are identical. Since these structure
constants are enough to reconstruct (locally) the internal product law of the
group, there should exist a strong link between the groups. This is made
precise by the following theorem:
THEOREM 18.30 The group SO(3) is a representation of SU(2): there exists a
homomorphism
R : SU(2) SO(3),
U 7 R U .
Moreover, for any U SU(2), the matrices U and U have same image: R U =
RU .
Proof. We need a way of constructing a matrix RU SO(3) for any matrix U
SU(2), such that the group structure is preserved. In other words, we need that
RU R V = RU V
for all U , V SU(2).
A trick is needed to do this. Denote by

M = M M2 (C) ; M = M , tr M = 0}
the vector space of hermitian matrices with trace zero. This may be written
z xi y
M =
; x, y, z R = { x ; x R3 }.
x+i y z
Thus, for any vector x = (x, y, z) R3 , we have a unique associated matrix in M ,

denoted M ( x) = x .
500
LEMMA 18.31 For all x, y R3 , we have
x x =

1
tr M ( x) M ( x )
2
det M ( x) = x 2 .
and
Conversely, any matrix M M can be written M =
xi i with xi = 21 tr (M i ).
Proof of the lemma. This is a simple computation.

Now we can describe the construction of the map U 7 RU .
Let U SU(2). We look for a matrix RU , that is, a rotation matrix, which we may
characterize by its action on R3 . Since M R3 , we may start by finding a matrix
acting on M . For this purpose, define
Q (U ) : M M ,
M 7 U M U 1 .
First notice that obviously we have Q (U ) = Q (U ) for all U SU(2). It is also

simple to check that M = U M U 1 is also in M : indeed, we have
M = (U M U 1 ) = (U M U ) = M
and
tr M = tr (U M U 1 ) = tr M = 0.
Hence there exists a unique element x R3 such that M = x . We then define
x = RU ( x), and in this manner we have defined a map RU : R3 R3 . The
uniqueness of x implies that RU is linear.
LEMMA 18.32 R U is a rotation.
Proof of the lemma. We compute the norm of x : we have

x 2 = det M = det(U M U 1 ) = det M = x 2 ,
and hence RU is an isometry. It is therefore in O(3). We may compute the
coefficients of RU , remarking that the vector x = (xi )i is mapped to x = (xi )i
with
P
xi = 21 tr (M i ) = 21 tr (U M U 1 i ) = 21 tr (U x j j U 1 i )
j
so
1
2
P
j
tr (U j U 1 i ) x j ,
1
tr (i U j U 1 ).
2
A direct computation shows finally that we have det RU > 0; hence RU is a
rotation.
(RU )i j =
To conclude, we check now that for any M M , we have

Q (U ) Q (V ) (M ) = U (V M V 1 )U 1 = (U V ) M (U V )1 = Q (U V )(M ),
showing that Q (U ) Q (V ) = Q (U V ) for any U , V SU(2) and therefore that U 7

Q (U ) is a homomorphism. It also satisfies Q (U ) = Q (U ). This implies that
these properties are also valid for the map U 7 RU (this may be checked using the
coefficients (RU )i j computed above).
Here is a summary of the course of the proof:

R
=n
rotation
vector
( defined up to 2)
M (
) =
matrix in M
= 21 tr (M
)
501
U = exp (i
/2)
matrix in SU(2)
( defined up to 4)
In other words, any element of SU(2) may be mapped to an element

in SO(3). However, this representation is not faithful: different elements
of SU(2) can have the same image in SO(3); indeed this is the case for U and
U which have the same associated representation (this is the only case where
this happens). In other words, there is no inverse to this homomorphism (it
would be bi-valued).
In fact, there is no continuous isomorphism between SO(3) and SO(2)
because these two spaces are not topologically equivalent:
THEOREM 18.33 SO(3) is not simply connected, but is doubly connected (i.e.,
running twice over any closed path gives a path which is homotopic to a point).
On the other hand SU(2) is simply connected.
Proof (intuitive). We will explain why SO(3) is not simply connected. Any element
R SO(3) may be characterized by a unit vector n, defining the axis of rotation, and
by its angle , as before, and R is parameterized again by the vector n. This is a
continuous parameterization, so the group SO(3) may be represented by a sphere with
radius (any rotation of angle between and 2 is also a rotation with angle [0, ]
and vector n).
This is an exact description, up to one small detail: the rotations ( n, ) and ( n, )
are equivalent. The group SO(3) is therefore, topologically, the same as the following
space: a sphere of radius on the surface of which diametrically opposite points have been
identified.
Hence the figure below represents a closed path in SO(3). It is indeed closed since
the north and south poles of the sphere are identified, that is, they represent the same
point in SO(3). The reader will easily convince herself that this closed path cannot be
continuously deformed into a single point (if we try to move the north pole toward
the south pole, the south pole is obliged to also move to remain opposite).
On the other hand, suppose we travel along this path twice. Then the resulting path
may be contracted to a single point, following the sequence of movements below (in
the first figure, the path taken twice is also drawn twice for legibility):
502

Note, moreover, that a closed path in SU(2) may be represented as a closed path
that is repeated twice in SO(3).
A consequence of this theorem is that a rotation of 4 is always equivalent

with a trivial rotation for a physical object.
What is this useful for?
You may perform the following magic trick. Take any object (a rectangular
pencil eraser, for instance), and attach at least two pieces of string to it. Then
attach the pieces of string to the walls, floor, or ceiling of the room where you
are operating. Then turn the eraser once (a 2-rotation). The strings become
entangled. Without cutting a string or turning the eraser again, it is impossible to
disentangle them and bring them back to their original state. However, if you turn
the eraser once more (another 2-rotation, so altogether a 4-rotation), it will
now be perfectly possible in front of an admiring and disbelieving audience
to disentangle the strings without any more turns or cutting (see sidebar 7,
page 506).
This shows that rotating the eraser twice (4-rotation) is equivalent to no
rotation at all, whereas one turn (2-rotation) is not.
But shouldnt a 2-rotation be, in fact, equivalent with a zero rotation, as is
commonly stated and believed? It all depends on the object considered. If, as
in this trick, the object is not by itself, the answer is No. More generally, the
mathematical properties of the object will reflect the path followed in SO(3)
to perform the rotation. In other words, the laws of the transformation must
reflect the group structure involved; but such a law is not unique: a physical
object may follow the trivial representation of SO(3) (it is then a scalar, such
as temperature), or a vector representation (it is then a vector, such as velocity);
a complex object may follow a representation of SU(2).
This is, for instance, the case of the wave function of the electron.
An electron may be represented (in nonrelativistic quantum mechanics) by
a wave function with two components,

( r, t)
( r , t) =
.
( r , t)
During a rotation, this changes according to a spinorial transformation:

( r , t)
( r, t)
/2
( r , t) 7 ( r , t) =
=e
,
( r , t)
( r , t)
(18.3)
where = n is the vector characterizing the rotation. Note that, as seen

previously, we have
but
e i(+2) n /2 = e i n /2 ,
e i(+4) n /2 = +e i n /2 .
Hence, a 4-rotation is required in order that the wave function remain invariant. How is it, then, that this particular behavior of the electron (it is not
Spin and Riemann sphere
503
invariant under a 2-rotation) has no consequence in physical observations?

After all, electrons exist abundantly in everyday matter. Note that, in fact, after a 2-rotation, only the sign of the wave function has changed; consequently,
the physical properties of the particle are unchanged. What changes are properties of interactions of the particles: in particular, the fact that electrons
are fermions (and hence obey the Pauli exclusion principle, so that matter is
stable) derive from this fact.
An object transforming according to the rule (18.3) is called a spinor.
Exercise 18.3 shows that other transformation rules are possible.
Remark 18.34 The fact that SU(2) is like twice SO(3) (or, more technically, it is a double
cover of SO(3)) is used to show that two states of spin exist in the fundamental spinor representation. In special relativity, the group of isometries of space is not SO(3) anymore, but
rather the group SO(3, 1), that leaves Minkowski space (R4 with the Minkowski metric) invariant. One can show that there exists a simply connected double cover of this group, isomorphic
to SL(2, C) and generated by two independent representations of SU(2) (one sometimes find
the notation SU(2) SU(2) to recall this independance). This is constructed using suitable
linear combinations of generators of rotations and Lorentz boosts as infinitesimal generators.
This group then affords four degrees of freedom for a particle with spin 21 , instead of the two
expected; this allows us to predict that not only the particle is thus described, but so is the
associated antiparticle.
The reader interested by the mathematics underlying the representations of the Dirac equation may begin by looking at the excellent book [68].
18.5
Spin and Riemann sphere
While we are dealing with pretty pictures, I will not resist the temptation
to present an application to quantum mechanics of the representation of the
complex plane as the Riemann sphere (see page 146); more precisely, this
concerns the representation of the spin.
We know that the wave function of an electron (or indeed of any other
particle with spin 12 ) may be represented as the tensor product of a vector | f
corresponding to a space dependency and a vector | corresponding to the
spin:
| = | f | .
This is another way of expressing that the wave function has two components,
as in (18.3) (the first, , corresponding to the state | or spin-up, and the
second, , corresponding to the state | or spin-down).
Let us concentrate only on the spin for now. Then the wave function |
is a linear combination of the basis vectors | and | . This, however, is
related to the choice of one particular direction for the measure of the spin,
namely, the direction z. We might as well have chosen another direction, for
504
instance, the direction x. Any vector | would then be a linear combination

of | and |, eigenvectors for the measure of spin on the x axis.
How is a wave function (denoted |) characterized when it is an eigenvector for
of spin in an arbitrarily chosen direction? Since the
the measure

family | , | is a basis of the space of spins, there exist complex numbers , C such that
| = | + | .
Because the physical meaning of a wave function is not changed when the
vector is multiplied by an arbitrary nonzero complex number, it follows that
only the ratio / is necessary to characterize the state |.
Then we have the following very nice result:
THEOREM 18.35 Let z = /; then we have z C = C {}. The stereo-
graphic projection z of z on the Riemann sphere gives the direction characterizing the
vector |.
z
1
1
i
z = /
For instance, to obtain the direction of positive x, we should take z = 1,

which gives (after normalizing appropriately):
i
1 h
| x = p | + | .
2
To obtain the direction (O x), we must take z = 1:
i
1 h
| x = p | | .
2
For the direction (O y), we take z = i:
i
1 h
| y = p | + i | .
2
Similarly, for the direction (O y), we take z = i:
i
1 h
| y = p | i | .
2
Exercises
505
Finally, for direction (O z), we must of course take z = :

| = | + 0 | ,
whereas the direction (O z) corresponds to z = 0, so:

| = 0 | + | .
EXERCISES
Exercise 18.2 Let i (i = 1, . . . , 3) denote the Pauli matrices. For a real number , compute
exp (i3 /2) k exp (i3 /2) when k = 1, 2, 3. Interpret this result.
Exercise 18.3 (Rotations acting on functions) Let R be a rotation.

To a smooth function
( r), we associate the function R such that R( r) = R 1 ( r) (so that the axes have been
transformed by the rotation R). For a rotation R z () around the O z-axis, show that we have
Rz ()(x, y, z) = (x y , y + x , z) and then that

dRz ()
= i L z Rz ()
with L z = i x
+y
.
d
y
x
Deduce that Rz () = exp(iL z ) and generalize this to the other axes. Show that
[i L i , i L j ] = i jk i Lk .
506
Sidebar 7 (Double connectedness of S O(3) and magic tricks) Start with a

sphere attached to the ceiling and floor with three strings (or more). At the beginning
of the experiment, the sphere is rotated by an angle of 4 (i.e., turned twice completely);
hence the strings that attach it become entangled. However, without turning the sphere
anymore, and of course without putting knife or scissors to the strings, it is possible to free
the strings, following the sequence below, at the end of which all the strings are exactly as
they where before the rotation was performed.
(It is necessary to assume that the strings are somewhat flexible; all this is topological,
that is, it requires continuous deformation.)
Initial positon.
After two rotations of the sphere, the

strings are entangled.
Beginning of the trick.
The white string goes over the rest.
The gray string starts moving.
Exercises
507
Sidebar 7 (Cont.)
The gray string goes behind the rest...
... continues...
... and goes back to its initial position.
The white string goes behind the black

string...
... and into position.
The black string is now free to take its

original place.
This magic trick (which may very well be done at home, using a pencil eraser, some tacks and
a yard or so of ordinary string) is a consequence of the double connectedness of SO(3).
Another illustration of this fact is Feynmanns hand trick [36].
Chapter
19
Introduction to probability
theory
Statistics is that science which proves
that 99.99% of human beings have
more legs than the average number.
The goal of probability theory is to provide mathematical tools to discuss rigorously

the elusive concept of randomness.
Randomness means in general that we have incomplete knowledge of some system.
Probability theory is interested, for instance, in what can be said of an event when
an experiment is repeated a large number of times. For instance, this experiment
may be the throwing of a die, and the event might be the result of the throw is at
most 3.
However, probability theory is not only concerned with random events, since it
is also at the heart of the formalism of statistical mechanics, which, at the classical
level, is based on deterministic equations (the Newton equations, in the Lagrangian
or Hamiltonian formalism). It can also be used to describe in a simpler (maybe
incomplete) manner some very complicated distributions, such as in signal theory
or in image processing.
Finally, it should be mentioned that probability theory may be used to discuss various events, and as the case may be to mislead, bluff, or hoodwink the unsuspecting
public. Although this is not something that students should do, it is important
that they know some of those techniques, not only to be capable of clear scientific
judgment (some experimental results are presented as proved,1 whereas rigorous
statistical analysis shows them to be no better than purely random), but also to
enlighten the daily life of an informed citizen.
510
Introduction to probability theory
19.1
Introduction
Dice are the embodiment of randomness. The Latin name, alea, is indeed
the source of the French word alatoire for random. The Arabic word
az-zahr, which also means dice, has produced Spanish azar, French hasard
(randomness), as well as English hazard. The way dice fall, through the
Latin cadere (to fall), has brought the Old French word chaance and then the
word chance.
According to a dictionary, chance is The unknown and unpredictable element
in happenings that seems to have no assignable cause.
This is what happens with the throw of a die. Even if we hold that the
laws that dictate the fall and rebounds of the die are perfectly deterministic,
it is a fact that this system is chaotic and thus, in practice, has unpredictable
behavior: the final outcome of the throw has no apparent reason [86].
It is quite remarkable that such a common concept as chance should at
the same time be so difficult to fathom. Here is a mathematical example:
Some real numbers have the following property,2 which we call Property P :
the decimal expansion of x contains statistically as many ze

ros as ones, twos, threes, etc; morevoer, any finite sequence of
digits (for instance, 123456789 or your birthdate) appears

P
in the decimal expansion and indeed appears statistically as
frequently as any other sequence of the same length (thus

12345 is as frequent as 94281).
Mathematicians (in particular, . Borel) have proved that Property P holds

for the overwhelming majority of real numbers. More precisely, the probability
that P holds for a real number x is equal to 1. The probability that P does
not hold is therefore equal to 0.
Now, can you give me a single example of a real number such that Property P
holds?
It should be noted that rational numbers do not have Property P . Indeed,
the decimal expansion of a rational number is periodic (for example, 1/7 =
0.142857142857... with an infinite repetition of the sequence 142857),
p which
clearly contradicts Property P . Maybe a number such as , or e, or 2 would
work. This is quite possible, but this is a complete mystery at the current
time: nobody has been able to show, for instance, that has Property P , nor
indeed that it doesnt.
So, we know that a randomly chosen real number must almost surely
satisfy Property P ; yet, (almost) any time we chose a number, either it does
1
2
E.g., the famous case of water memory, which is described in detail in [15].
A number for which Property P holds is called a normal number.
Introduction
511
not have this property, or if it does, we are not able to prove it.3 The conclusion
is that it is in fact very difficult to make a random choice of a real number;
most of our choices are not random at all.4
This simple example should be enough to suggest how difficult it may be to
express rigorously what is a random choice. To overcome those difficulties,
it has been necessary to formalize and axiomatize probability theory, leaving
aside questions of interpretation to concentrate solely on the mathematical
content.
In this sense, the forefathers of probability are certainly Blaise Pascal (see
page 606) and Pierre de Fermat (16011665), during a famous exchange of letters during the summer 1654. The physicist Christiaan Huyghens (16291695)
published this correspondance in the first treatise of probability theory, De Ratiociniis in ludo ale, where in particular the notion of mathematical expectation
is introduced.
However, the works of Jacques Bernoulli (see page 555) and Abraham de
Moivre,5 with the first theorems concerning laws of large numbers, are the
real starting point of modern probability theory. Laplace and Gauss studied
the theory of errors (with important applications to measurements such as
that of the terrestrial meridian).
In the ninetheenth century, Pafnouti Tchebychev6 (18211894) and his
students Andrei Markov (18561922) and Alexandre Liapounov (18571918)
created a formalism suitable for the study of sums of random variables,
which marks the beginning of the important Russian school of probability
theory. The complete axiomatization of probability is due to Andrei Kolmogorov (see below), building on previous work of Sergei Bernstein (1880
1968), Richard von Mises (18831953) and mile Borel (18711956).
3
To be precise, we know a few normal numbers. For instance, the number obtained by first
writing successively all positive integers, then concatenating the resulting sequences of digits
(Champernowne, 1933):
0, 1234567891011121314151617 . . .
or the number obtained by doing the same with the sequence of prime numbers (Erds and
Copeland, 1945) :
0, 23571113171923293137 . . .
It is quite obvious that those numbers are artificial; in particular, they depend on the fact
that we work in base 10. So one may increase the difficulty by asking for an absolutely normal
number, one which is normal (in an obvious sense) in any basis. It seems that today not a
single absolutely normal number is known, whereas, again, almost all numbers must have this
property.
4
Consider also the great difficulty inherent in the writing of a random number generator
for a computer.
5
Abraham de Moivre (16671754) left France at eighteen during the repression against
Protestants. He discovered mathematics through a chance encounter with Newtons Principia [69]. Besides studies of complex numbers, he gave the mathematical definition of independent events, and proved the Stirling formula.
6
The mathematical great-great-great-great-grandfather of the translator.
512
19.2
Basic definitions
In this section, we introduce the basic vocabulary of probability theory
and the axioms stated by Kolmogorov [19] in ber die analytischen Methoden in der Wahrscheinlichkeitrechnung (Analytic Methods of Probability Theory,
1931), and also in the historical book Grundbegriffe der Wahrscheinlichkeitrechnung
(Foundations of Probability Theory, 1933).
The probability of an event that can range only over a discrete set (such as
the result of throwing a die) can be captured very intuitively. However, when
the space of possible events is continuous (such as the points in space, or an
instant in time), a precise mathematical description of the probability that
an event occurs is required.
Kolmogorov suggested that the probability of an event be the expression
of the measure of a set.
DEFINITION 19.1 A probability space is a pair (, ), where
: a set (the sample space).

: -algebra of events, that is, a set such that
i)
ii)
iii)
iv)
P() (P() is the set of all subsets of );

and ;
if A , then its complement A also belongs to ;
S
if An for any n N, then nN An .
An element of is called an event. An element of is sometimes called

the result of the experiment , or an atomic or elementary event.
The reader should recognize Definitions 2.7 and 2.10 on page 59. So a
measurable space is a probabilistic space.
Example 19.2 When playing (once) heads or tails, one can take = {P , F } and

= , {P }, {F }, {P , F } .
Example 19.3 When rolling once a die with six facets, the space of elementary events can be
= {1, 2, 3, 4, 5, 6}
and
= P().
The event the result of the throw is even is the element {2, 4, 6} .

Example 19.4 Take = R. Then both {, R} and , [0, 1] , R [0, 1] , R are -algebras.
Remark 19.5 Note that it is sometimes useful to take for a set which is not an abstract
mathematical object, but rather a group of persons, a collection of objects, etc.
Basic definitions
513
DEFINITION 19.6 An event A is realized if the result of an experi-
ment is in A.
Working with the -algebra , we have at our disposal on the set of

events the usual set-theoretic operations of union, intersection, and taking the
complementary subset. These have the following intuitive meaning:
the union of two events A and B, that is, A B, corresponds to the
realization of A or B (an inclusive or, where A and B is possible);
the intersection of two events A and B, i.e., A B, corresponds to the
realization of A and B;
taking the complementary subset corresponds to the realization of the
opposite of A, or equivalently to the no-realization (the failure) of A.
DEFINITION 19.7 Let (, ) be a sample probability space A , B . Two
events A and B are simultaneously realized if the outcome of the experiment
is A B.
The events A and B are incompatible if A B = .
A complete class of events is any countable (or finite) sequence of nonempty events (An )nN , which are mutually incompatible and which together
exhaust , that is, such that
i) Am An = if m 6= n,
S
ii) nN An = .
Any elementary event {} P() is then contained in one and only one of
the elements of the complete class.
DEFINITION 19.8 Let C P() be any set of subsets of . The intersection
of all -algebras containing C is a -algebra, which is called the -algebra

generated by C and is denoted (C
C ).
Example 19.9 Consider the experiment of rolling a die once. We take
= {1, 2, 3, 4, 5, 6} and we put = P().

Consider now C = {1, 2} . The -algebra generated by C is

(C ) = , {1, 2}, {3, 4, 5, 6}, .

On the other hand, if we put C = {1}, {2} , we have

(C ) = , {1}, {2}, {1, 2}, {3, 4, 5, 6}, {1, 3, 4, 5, 6}, {2, 3, 4, 5, 6}, .

Having defined these notions, which are basically set-theoretic in nature,

we can use the -algebra for the purpose of defining measures of sets in this
-algebra, as in the definition of the Lebesgue integral. A suitably normalized
positive measure will provide the notion of probability.
514
DEFINITION 19.10 (Probability measure) Let (, ) denote a sample space
with a -algebra. A map P : R+ is a probability (measure) on (, ) if:

(P1) P() = 1;
(P2) for any countable (or finite) sequence (An )nN of pairwise disjoint
events in , we have
S P
An =
P(An ).
P
nN
nN
The triplet (, , P) is also called a probability space.
Comparing with Definition 2.11 on page 59, we see that a probability space
is simply a measure space where the total measure of the space is supposed to
be equal to 1.
The reader is then invited to check the following properties:
PROPOSITION 19.11 Let P be a probability measure. Then the following properties
hold:
i) if A B, then we have P(A) P(B);

ii) P A = 1 P(A);
iii) if A B, then we have P(B) = P(A) + P(B \ A);

iv) P(A B) = P(A) + P(B) P(A B);
v) if (An )nN is an increasing sequence (for inclusion), then we have

S
P
An = sup P(An ) = lim P(An );
nN
nN
vi) if (An )nN is a decreasing sequence (for inclusion), then we have

T
P
An = inf P(An ) = lim P(An ).
nN
nN
Example 19.12 Let a be an elementary event. The map
a : P() {0, 1},
1 if a A ,
A 7
0 if a
/A
is a probability measure, which is called the Dirac measure or mass at a.

Example
19.13 Let (pn )nN be an arbitrary sequence of non-negative real numbers such that
P
p = 1. Let be any set and let (xn )nN be a sequence of elements
in . Then

PnN n
p
defines
(in
an
obvious
way)
a
probability
measure
on
,
P
()
.
nN n xn
Basic definitions
515
Andrei Nikolaievitch Kolmogorov (19031987) could not receive any

education from his parents, as his mother died in childbirth and his
father, an agriculturist, was in exile. He was raised by his aunt and
became a train conductor. In his free time, Kolmogorov wrote a
treatise on Newtonian mechanics. His further works dealt with very
varied subjects: trigonometric series, the axiomatization of probability theory, the theory of Markov processes, information theory. In
physics, his contributions were important in the theory of turbulence, the study of chaotic systems and dynamical systems (his initial
remains in the acronym KAM theory, with Arnolds and Mosers).
He also had deep interests outside mathematics, including the poetic
works of Pushkin.
Example 19.14
Let f : Rn R+ be an integrable function (with respect to Lebesgue measure)
R
such that f d = 1. For any Borel subset A B(Rn ) (recall that the Borel -algebra is the
-algebra on Rn generated by open sets, or equivalently by rectangles [a1 , b 1 ] [an , b n ]),
let
Z
f d.
P(A) =

The map P thus defined is a probability measure on Rn , B(Rn ) .
Example 19.15 Let be a positive real number, and let m R. The Gauss distribution or
normal distribution with mean m and standard deviation is the probability measure on R
(with respect to the Borel -algebra) defined by

Z

(t m)2
1
dt
for any A B(R).
P(A) = p
exp
2 2
2 A
DEFINITION 19.16 A subset N is negligible (for a given probability

measure P) if it is contained in a measurable subset (and event) with probability zero:
N
N N and P(N ) = 0.
A property p that depends on elementary events holds almost surely or
is almost sure if the set of those which do not satisfy p, namely, the set

; p() is false ,
is negligible.
Example 19.17 Let = [0, 1], with the Borel -algebra = B [0, 1]
and the uniform

probability measure on [0, 1] (i.e., the Lebesgue measure restricted to [0, 1]). Then for x [0, 1],
x is almost surely an irrational number.
516
A
B
C
Fig. 19.1 The Poincar formula for n = 3. In dark gray, A B C ; in light gray, A B,
A C and B C .
19.3
Poincare formula
Consider two events A and B. Obviously (just using the definition as the
measure of the corresponding events) we have the formula
P(A B) = P(A) + P(B) P(A B).
Can this be generalized to more than two events? Let A, B, C be three

events. We can compute
P(A) + P(B) + P(C ) P(A B) P(B C ) P(C A),
but as clearly seen in Figure 19.1, this is not the correct result for P(A B C ),
because the subset A B C has been counted three times, then subtracted
three times. Hence the correct relation is in fact
P(A B C ) = P(A) + P(B) + P(C ) P(A B)
P(B C ) P(C A) + P(A B C ).
With this done, the general result is easy to guess and prove.
THEOREM 19.18 (Poincare formula) Let n N and let A1 , . . . , An be arbitrary
events, not necessarily disjoint. Then we have

n

n
S
P
P
P
Ai =
(1)k+1
i=1
k=1
1i1 <i2 <<ik n
P(Ai1 Aik ).
Proof. A proof using random variables, as defined in the next chapter, is given in
the Appendix D. An elementary proof by induction is also easy.
Conditional probability
517
19.4
Conditional probability
Exercise 19.1 Youre taking a plane to Seattle. A colleague tells you: There is a one in ten
thousand chance that there is a bomb in a plane. But there is only one in a hundred million
chance that there are two bombs. So, for more security, just bring your own bomb.
What should you think of this argument?
DEFINITION 19.19 Let A and B be two events such that P(B) 6= 0. The
conditional probability of the event A knowing the event B is the quantity

P B (A) = P(A|B) =
P(A B)
.
P(B)
It is denoted either P B (A) or P(A|B).

Example 19.20 The probability that there are two bombs in your plane, knowing that there is
at least one (yours) is therefore simply one in ten thousand. Bringing one in your luggage is a
useless precaution.
If P(A) 6= 0 as well as P(B) 6= 0, it is also possible to define the probability

of B knowing A. Then we have
P(A B) = P(B) P(A|B) = P(A) P(B|A),
hence the following formula, due to the clergyman Thomas Bayes, who published it in 1764 [12]:
THEOREM 19.21 (Bayes formula) Let A and B be two events in with nonzero
probability. Then we have
P(B|A) =
P(A|B) P(B)
.
P(A)
Although it is very simple, this is an interesting result since it relates

the probability that the event B is realized knowing A to the reciprocal
probability that A be realized knowing B. One has to be careful because
those two quantities are often confused, either knowingly or unknowingly.
For instance, a scary message such as 95% of heroin addicts started by
smoking marijuana may be used with the intent of suggesting (without stating
it explicitly) that 95% of smokers of marijuana become heroin addicts. This
is the ordinary fare of most political campaign arguments, of advertisements
of every kind, and those never-ending polls.
But the greatest care is necessary, as the following example shows:
518
Thomas Bayes (17021761), English Presbyterian minister and theologian, studied mathematics, and in particular probability theory,
during his free time. He published little of his work during his
lifetime and remained unknown to most of his peers. His Essay
towards Solving a Problem in the Doctrine of Chances was published
posthumously in 1764.
Example 19.22 On the planet Zork live two ethnic groups: Ents and Vogons. According to a
recent statistical analysis of the distribution of wealth on Zork, it is true that
80% among the Vogons are poor
80% among the poor are Vogons.
Is it possible to deduce that wealth is inequitably distributed between the two groups?
Solution: Certainly not, in the absence of other data such as the proportion
of Vogons in the population. It may very well be the case that the Vogons
represent 80% of the total population, and that 80% of the total population is
poor. If that is the case, wealth is equitably distributed between the two groups
(not necessarily within each group). Again, this kind of confusion occurs
continually, and is sometimes encouraged. (Many examples can be found by
reading almost any newspaper, adapting the two words poor and Vogons
to various circumstances.)
We now introduce a minor generalization of a complete system, which

allows the events in the system only to cover the sample space up to a set of
probability zero.
DEFINITION 19.23 An almost complete system is a sequence (An )nN of
events in such that
i) Am An = if m 6= n,
S

ii) P
A
= 1.
n
nN
THEOREM 19.24 (Bayes formula) Let (An )nN be an almost complete system of
events, and let B be an arbitrary event B . Assume that P(An ) > 0 for any
n N. Then we have
P
i) P(B) = nN P(B|An )P(An );
ii) if, moreover, P(B) > 0, then
p N
P(B|A p ) P(A p )
P(A p |B) = P
.
P(B|An ) P(An )
nN
Independent events
519
Remark 19.25 The Collins case is a striking example of the use (or misuse) of probability and
conditional probability. In Los Angeles in 1964, a blond woman and a black man with a beard
were arrested for a robbery. Despite the absence of convincing evidence, the prosecution argued
successfully that they must be guilty because the chance that a random couple corresponded (as
they did) to the witnesss description was estimated to be one in twelve million. The California
Supreme Court reversed the judgment on appeal, since it was shown that the probability that
at least two couples in the Los Angeles area correspond to the description knowing that at least
one couple does (namely, the actual thieves) was close to 42 %, and therefore far from negligible
(certainly too large to decide guilt beyond reasonable doubt!). This shows that there was a
high probability that arresting the couple was a mistake.
19.5
Independent events
It is now time to define the notion of independent events. Intuitively, to
say that two events A and B are independent means that the probability that
A is realized is equal to the probability that it is realized knowing that B is
realized, or indeed to the probability that A is realized knowing that not-B
is realized.
DEFINITION 19.26 Let A and B be two events in . Then A and B are
independent if P(A B) = P(A) P(B), or, equivalently if 0 < P(B) < 1, if

P(A|B) = P(A|B) = P(A).
Let 1 , . . . , n be -algebras contained in . They are independent algebras if
T
Q
n
n
(B1 , . . . , Bn ) 1 n ,
P
Bi =
P(B i ).
i=1
i=1
Moreover, events A1 , . . . , A n are independent if the -algebras

are independent.
({A1 }), . . . , ({An })
Warning! The pairwise independence of events A1 ,. . . , An does not imply

their independence. In fact, we have the following result:
THEOREM 19.27 The events A1 , . . . , An are independent if and only if, for any
p N, p n, and for any p-uple of distinct indices (i1 , . . . , i p ), we have

T
Q
p
p
P
A ik =
P(Aik ).
k=1
k=1
In particular, for n = 2, we recover the first definition: A and B are independent

events if and only if P(A B) = P(A) P(B).
520
Example 19.28 Consider the experiment of rolling two dice, and the events:
event A: the sum of the results of both dice is even;

event B: the result of the first die is even;
event C: the result of the second die is even.
It is easy to check that A, B, and C are pairwise independent. Indeed, we have (assuming all
36 outcomes of the experiment have probability 1/36):
P(A) = P(B) = P(C ) =
1
2
whereas P(A B) = P(B C ) = P(C A) = 41 .
On other hand, A, B, C are not independent, since

P(A B C ) = P(B C ) =
1
4
whereas P(A) P(B) P(C ) = 81 .
(This is also clear intuitively, since A and B together imply C , for instance.)
One should note that this concept of independence depends only on a

numerical evaluation of probabilities, and that it may turn out to be in contradiction with common sense or intuition, for which independence may carry
deeper meaning. In particular, independence should not be mistaken with
causality (see Section 20.11, page 550).
Chapter
20
Random variables
Falstaff. Prithee, no more prattling; go. Ill hold.
This is the third time; I hope good luck lies in odd numbers.
Away I go. They say there is divinity in odd numbers,
either in nativity, chance, or death. Away!
William Shakespeare
The Merry Wives of Windsor, Act V, Scene i.
As we have noticed in the previous chapter, the sample space is not

usually a mathematically simple set it may well be a population of exotic
fishes, a collection of ideas, or representatives of various physical measurements. Worse, this set is often not known. Then, instead of working with
the probability measure P on , one tries to reduce to various probabilities
on R.
For this purpose, we consider a map X from to R (which is called
a random variable); R becomes the new set on which a measure is defined.
Indeed, a probability measure denoted P X , which is the image of P by X ,
and which is called the probability distribution of X , is defined.
20.1
Random variables and probability distributions
DEFINITION 20.1 Let (, ) be a probability space. A random variable is a
function X , defined on and with values in R, which is measurable when R
Random variables
522
is equipped with the Borel -algebra, that is, such that we have X 1 (A)
for any Borel set A B(R).
To simplify notation, we will abbreviate r.v. instead of writing random
variable.
A random function is, in other words, simply a measurable function
and, for a physicist not overly concerned with quoting the axiom of choice,
this means any function that she may think about.
Here is an example. Take as probability space
def
= {all students at Rutgers university},

with the -algebra = P().
We can define a random variable X corresponding to the age of a student:
X (student) = age of the student R.
Now we can study the distribution of the age in the population (evaluate
the average, the standard deviation,...) without knowing too much about ,
which can be a rather complicated set to envision precisely (if we think that
it contains information such as the fact that Franois prefers to sing Bach in
B-flat, or that Winnies preferred dessert is strawberry cheesecake).
Given a probability space (, , P) and an r.v. X on , the probability
distribution of X is defined as the image probability measure of P by X :
DEFINITION 20.2 (Distribution) Let (, , P) be a probability space and X
a random variable. The probability distribution of X is the probability
measure
P X : B R+ ,
B 7 P{X B},
where

def
{X B} = X 1 (B) = ; X () B .
In particular, we have

P X ]a, b] = P{a < X b} = P X 1 ]a, b]

= P ; a < X () b .
Note that this probability measure is in general defined on another set

(R) and with respect to another -algebra (of Borel sets) than that of the
probability space itself ().
The following diagram summarizes this discussion:
Random variables and probability distributions

(unknown set)
(unknown)
P
?
523
X R
P X (known)
?
In the example of the Rutgers students, assume moreover that the sample
space is given the uniform probability measure (each student is given the same
probability, p = 1/N , where N is the total number of students). What is the
probability distribution of the random variable age, which we see as integervalued? This distribution is a probability measure on the set N, an instance
of what is often called a discrete distribution:

P X {n} = P {X = n} = P {students ; age (student) = n}
number of students aged n
.
N
It can very well be the case that two random variables have the same
distribution without being closely related to each other. For instance,
=
THEOREM 20.3 Let X and Y be two random variables defined on the same sample
space which are equal almost everywhere. Then X and Y have the same distribution:
PX = PY .
But be warned that the converse statement is completely false.

Consider, for example, the process of flipping n times a fair (unbiased)
coin, using 0 and 1 to denote, respectively, heads or tails. We can take
= {0, 1}n and of course = P(). Let = (E1 , . . . , En ) , where
E i {0, 1} for i = 1, . . . , n, be an atomic event. The probability of {} ,
if the successive throws are independent, is equal to

P {} = (1/2)n
(uniform probability).
Consider now the random variable
X : N,
n
X
7
E i = number of tails.
i=1
We can find the distribution of X by denoting, for instance, Ak the event

corresponding to k tails out of n throws. So Ak and in fact

Ak = ; X () = k .
The distribution of X is supported, in the case considered, on the set {0, . . . , n},
and we have
n

n
1
P X {k} = P(X = k) = P(Ak ) =
.
k
2
Random variables
524
This distribution
is called the binomial distribution, and it is often denoted

B n, 12 . Note that it is easily checked that

n
1 X
n
1
P X (N) = n
= n (1 + 1)n = 1,
2 k=0 k
2
as it should be. Moreover, it is easily seen that if X has a binomial distribution, so does (n X ) (corresponding to the number of heads in n throws).
Of course, X and (n X ) are not equal almost everywhere!
What is clearly visible from this example is that distributions of random
variables only provide some information on the full sample space, by no means
all. This is the price to pay for the simplification that results from their use.
Whatever is not encoded in the random variable (all the factors leading to a
student of a certain age and not another, the exact trajectory in space of the
falling die...) is completely forgotten.
DEFINITION 20.4 (Bernoulli distribution, binomial distribution) Let n N
and p [0, 1]. A random variable X follows the Bernouilli distribution

with parameter p, denoted B (1, p), if

and
P X {0} = 1 p.
P X {1} = p
A random variable X follows the binomial distribution B ( n, p) if

n k
P X {k} =
p (1 p)nk
for k [[1, n]].
k
The binomial distribution B(n, p) corresponds to the sum of n independent random variables, each of which has a Bernoulli distribution with the
same parameter B(1, p).
20.2
Distribution function and probability density
DEFINITION 20.5 The distribution function of an r.v. X is the function
defined for x 0 by

def
F (x) = P {X x} = P X ], x] .
For all a, b R, we have

P {a < X b} = F (b) F (a).
Distribution function and probability density
525
Since the -algebra of Borel subsets of the real line is generated by the
intervals ], a], the folowing result follows:
PROPOSITION 20.6 Let (, ) be a probability
space. A map X : R is a
random variable if and only if X 1 ], a] for all a R.
Moreover, the probability
distribution P X is uniquely determined by the values

P X 1 (], a]) = F (a) for a R, that is, by the distribution function.
All information contained in the probability distribution P X is contained

in the distribution function F . The distribution function is thus a fundamental tool of probability theory.
Exercise 20.1 Express P {a < X < b } , P {a X b } , and P {a X < b } in terms
of F (b ), F (a), F (b ), and F (a ).
From the definition of the distribution function, we immediately derive

the following theorem:
THEOREM 20.7 Let X be a random variable with distribution function F . Then
F has the following properties:
i) F is right-continuous; that is, for all x R, we have

def
F (x + ) = lim+ F (x + ) = F (x) ;
0
ii) F is not necessarily left-continuous, and more precisely we have

F (x) F (x ) = P{X = x} ;
iii) for all x R, we have 0 F (x) 1 ;
iv) F is increasing ;
v) if we define
def
def
F () = lim F (x)
and
F (+) = lim F (x),
F () = 0
and
F (+) = 1 ;
x+
then
vi) F is continuous if and only if P{X = x} = 0 for all x R. Such a
probability distribution or random variable is called diffuse.
Proof. Only i) and ii) require some care. First, write F (x + ) F (x) = P{x <
X x + }, then notice that, as tends to 0 (following an arbitrary sequence), the set
{x < X x + } tends to the empty set which has measure zero. This proves i).
Moreover, writing F (x) F (x ) = P{x < X x}, one notices similarly that
the set {x < X x} tends to {X = x}.
Random variables
526
20.2.a
Discrete random variables
DEFINITION 20.8 Let (, , P) be a probability space. A random variable X
on is a discrete random variable if X () is almost surely countable, that

is, there is a countable set D such that P{X
/ D } = 0.
When an r.v. is discrete, it is clear that its distribution function has

discontinuities; it is a step function.
Example 20.9 Consider the experience of throwing a weighted die, which falls on the numbers
1 to 6 with probabilities respectively equal to 14 , 16 , 16 , 16 , 16 , 121 (this is the one I hand to my
friend Claude when playing with him).
The distribution function of the corresponding random variable is given by the following
graph:
1
5/6
4/6
3/6
2/6
1/6
1
20.2.b (Absolutely) continuous random variables

DEFINITION 20.10 A random variable is continuous if its distribution func-
tion is continuous.
Remark 20.11 In fact, we should not speak simply of continuity, but rather of a similar but
stronger property called absolute continuity. The reader is invited to look into a probability
or analysis book such as [83] for more details.
Since, for distribution functions, continuity and absolute continuity are in fact equivalent,
we will continue writing continuous instead of absolutely continuous for simplicity.
DEFINITION 20.12 The probability density of a continuous r.v. with distri-
bution function F is a non-negative function f such that for all x R we

have
Zx

P {X x} = F (x) =
f (t) dt.
The density probability is almost everywhere equal to the derivative of F .
For a discrete r.v., the probability density f does not exist as a function, but
it may be represented with Dirac distributions, since we want to differentiate
a discontinuous function.
Expectation and variance
527
Example 20.13 Consider again the weighted die above. The probability density of the r.v. X
representing the result of throwing this die is
1
1
1
1
1
1
f (x) = (x 1) + (x 2) + (x 3) + (x 4) + (x 5) + (x 6).
4
6
6
6
6
12
Remark 20.14 In many cases, a random variable has a distribution function which can be
expressed as the sum of a continuous distribution function and a step function (which is itself
a distribution function). The probability density is then the sum of an ordinary function and
at most countably many Dirac distributions.
20.3
20.3.a Case of a discrete r.v.
DEFINITION 20.15 Let X be a discrete random variable. Assume that
X () = { pk ; k I } D ,
where I is a countable set and P{X D } = 0, and let pk = P(X = xk ) for

k I . The expectation of X is the quantity
X
X
E(X ) =
xk P(X = xk ) =
xk pk
kI
kI
if this series converges absolutely.
Thus, the expectation is the expected average of the random variable. We

believe that in some ways it represents the typical behavior of X . To make this
more precise, we try to characterize how the values of X can be far from the
expectation. For instance, during an exam graded on a 200-point scale, the
class average may be equal to 100 with grades all contained between 90 and
110, or with grades ranging from 0 to 200; although two such distributions
may yield the same expectation, they are of course very different.
To measure the deviation from the mean, the first idea might be to consider
the average of the differences from the expected value, that is, to consider
X

E X E(X ) =
xk E(x) pk ;
k
however, this is always zero, by the definition of the expectation. Another idea
is therefore needed, and since we want to add all the deviations, irrespective
of sign, we can try to use

X

xk E(x) pk .
E X E(X ) =
k
Random variables
528
However, absolute values are always algebraically delicate to handle. It is better,

therefore, to get rid of the signs by taking the square of X E(X ), possibly
taking the square root of the result obtained by taking the average of this.1
DEFINITION 20.16 The variance of the discrete r.v. X is the quantity

2 P
2
Var(X ) = 2 (X ) = E X E(X )
=
xk E(X ) pk
k

= E(X 2 ) E(X )2
which may be either a non-negative

p real number or +. The standard deviation of X is given by (X ) = Var(X ).
Consider the example of giving global grades to students having passed

four exams, say in Mathematics, Physics, Chemistry, and Medieval French
Poetry. All four courses should play a similar role in determining the ranking
of the students, even though they may carry different coefficients. For this
purpose the graders must ensure that the standard deviation of their own
grades is equal to some value fixed beforehand with their colleagues (say, 4). It
is then of no consequence that the average in Mathematics is different from
the average in Medieval French Poetry. On the other hand, if the standard
deviation in Mathematics is equal to 6 while that in Medieval French Poetry
is 2, then the Medieval French Poetry exam has little effect in distinguishing
the students (they all have more or less the same grade), whereas Mathematics
is crucial, even if this being the entrance exam to the French Academy
the coefficient for Mathematics is a factor of 2 smaller.
20.3.b Case of a continuous r.v.

We will take the following definition as generalizing the discrete case.
DEFINITION 20.17 Let X be a random variable with probability density f .
The expectation E(X) of the r.v. X is the number

Z +
E(X ) =
t f (t) dt
when the integral exists. The variance Var(X) and the standard deviation
(X) are given by
Z +
Z +
2

2
Var(X ) = 2 (X ) =
t E(X ) f (t) dt =
t 2 f (t) dt E(x)
when those integrals exist.

1
In order to preserve the unit of a physical quantity, for instance.
529
X is discrete, and a
distribution is used to represent its probability density, then the above definition provides the
same definition of the expectation of X as Definition 20.15.
Remark 20.18 The reader may check, as an exercise, that if an r.v.
2
Remark 20.19 Putting f (x) = (x) , the reader can also recognize in the definition of
expectation and standard deviation the quantities denoted x and x in quantum mechanics
(see Section 13.3 on page 359).
Remark 20.20 The expectation and the variance are special cases of what are called moments of
random variables.
For a continuous random variable, we write equivalently

Z
Z
x f (x) dx =
X () dP().
R
The first formula uses the probability density;2 we can consider that we
integrate the function x 7 x with respect to the measure d = f (x) dx.
The second formule uses an integral over the sample space , which
is equipped with a -algebra and a probability measure, which suffices
to defined an integral in the same manner that the Lebesgue integral is
constructed.3 The function considered is this time the random variable
X , and it is integrated with respect to the measure dP().
Mathematicians usually avoid using the density f (which requires the use
of the theory of distributions as soon as F is discontinuous) and appeal instead to another measure denoted dF = f (x) dx and to the so-called Stieltjes
integrals (see sidebar 8 on page 552). However, since density probabilities
are probably more intuitive to physicists, we will use those instead, it being
understood that it may be a distribution, and must be considered accordingly,
even if it is written under an integral sign.
Exercise 20.2 A random variable X is distributed according to a Cauchy distribution if its
probability density is of the type
f (x) =
a
1
a 2 + (x m)2
for some real number m and some positive real number a. Compute the expectation and
standard deviation of a random variable with this distribution (answer in the table, page 616).
2
And, as already noticed, may be generalized to the case of discontinuous F by considering
f as a distribution and writing f , x 7 x instead of the integral.
3
This is indeed what is called the Lebesgue integral with respect to an arbitrary measure on
an arbitrary measure space.
Random variables
530
20.4
An example: the Poisson distribution
20.4.a
Particles in a confined gas
Consider a box with constant volume V , containing N identical, classical

(i.e., not quantum) particles, which we assume not to interact. Moreover,
assume that the system has reached thermodynamic equilibrium, and is in the
situation of thermodynamic limit (in particular, N is supposed to be very large).
Fix now a small volume V V . What is the probability that one
given particle is located in this small volume? If the system is ergodic, this
probability is p = V /V .
What is now the probability that there are two particles in this volume?
That there are n particles? More generally, we ask for the probability distribution of the number of particles contained in V .
To start with, notice that since the particles are classical and do not interact,
there is no correlation between their positions. Thus, the probability for
one particle to be in the given volume V is independent of the number
of particles already located there. Hence, the probability distribution for the
number of particles in V is the binomial distribution P = B(N , p) with
values given by

N
P(n) =
p n (1 p) N n ,
n
since n particles must be in V (with probability p) and N n must be outside (with probability (1 p)). The expected value of the binomial distribution
is
N
X
n=
n P(n) = N p.
n=0
If N goes to infinity, we will see in the next chapter that the binomial distribution can be approached by a gaussian distribution; however, this is only
correct if N p 1. For fixed N (which is the case in physics), if we consider
a very small volume V , it is possible to make n = N p = N V /V very
small compared to 1. Then we can write

N
Nn
and
(1 p) N n = exp(N n) log(1 p) e N p = e n ,
n
n!
which gives
P(n)
N n n n nn n
p e = e .
n!
n!
An example: the Poisson distribution
531
DEFINITION 20.21 The Poisson distribution with parameter is the discrete
probability measure such that
n
for all n N.
n!
A random variable X is a Poisson random variable if P{X = n} = P (n) for
all n N.
P (n) = e
The Poisson distribution is useful to describe rare events.4 The maximal

value of n 7 P (n) is attained for n . Moreover, the reader will easily
check the following results:
THEOREM 20.22 The expectation of a Poisson random with parameter is E(X ) =
, and its variance is Var(X ) = .
20.4.b Radioactive decay

A radiation source randomly emits -particles. Assume that this radiation
is weak enough that it does not perturb the source (i.e., the power which is
radiated is statistically constant). Moreover, assume that the probability that
a particle is emitted between t and t + t depends only on t, and not on
t or on the number of -particles already emitted before time t (this is called
a Markov hypothesis). It seems reasonable to postulate that this probability
is proportional to t when t is very small; thus, we take it to be equal to
pt.
Let Q(t) be the probability that the source emits no particle during a
length of time of duration t. In the absence of time correlation, we have
therefore
Q(t + dt) = (1 p dt) Q(t),
which, together with the condition Q(0) = 1, implies that Q(t) = e p t .
What is the probability that n particles are emitted during an interval of
length t? To compute this, partition the interval [0, t] into a large number
N of small intervals of length t = t/N . During each of these intervals
of length t, the probability of emitting an -particle is pt. Hence the
probability of emitting n particles is given by a binomial distribution, which
may be approximated by the Poisson distribution with parameter = p t:

N
(p t)n p t
P(n, t) =
(pt)n (1 pt) N n
e .
n
n!
As might be expected, this probability is maximal when n p t; the
average number of particles emitted during an interval of length t is given by
the expectation of the Poisson distribution, and is also equal to p t.
4
The first recorded use of the Poisson distribution is found in a statistical study of the
annual number of deaths caused by horse kicks in the Prussian cavalry in 1898.
Random variables
532
20.5
Moments of a random variable
DEFINITION 20.23 The k-th moment (with k N) of a random variable X
with distribution function F is the quantity

Z +
Z
def
k
k =
x f (x) dx =
X k () dP()
if it exists.
For k 1, the k-th moment does not necessarily exist. The 0-th moment
always exists, but it is always equal to 1 and is not very interesting. The
first moment is simply the expectation of X . One can define variants of the
moment for a centered random variable, that is, one which expectation equal
to zero.
DEFINITION 20.24 A centered random variable is a random variable with
expectation equal to 0. If X is a random variable such that the expectation of

X exists, and m = E(X ), then (X m) is a centered random variable.
DEFINITION 20.25 The k-th centered moment of an r.v. X with expectation
m is given by
def
k =
(x
m)k
f (x) dx =
X () m)k dP().
Remark 20.26 For k 1, the centered k-th moment exists if and only if the original (noncentered) k-th moment exists.
Remark 20.27 The centered second moment is none other than the variance:
2 = Var(X ) = 2 (X ).
Exercise 20.3 Show that 2 =
2 21 .
Show that the kinetic energy of a distribution of mass is the sum of the kinetic energy
in the barycentric referential and the kinetic energy of the center of gravity (also called the
barycenter), carrying the total mass of the system.
The following definition is sometimes used.

DEFINITION 20.28 A reduced random variable is a random variable for
which the standard deviation exists and satisfies (X ) = 1.

If X is an r.v. with expectation m and standard deviation , the random
variable (X m)/ is always centered and reduced.
Moments of a random variable
533
DEFINITION 20.29 The median of a random variable X with distribution
function F is the number
1
.
x1/2 = inf x R ; F (x)
2
If F is continuous, then x1/2 is a solution of the equation F (x) = 12 . If,

moreover, F is strictly increasing, this is the unique solution.
Example 20.30 The median occurs frequently in economics and statistics. The median salary
of a population, for instance, is the salary such that half of the population earn less, and half
earns more. It is therefore different from the average salary.
Take for instance the case of six individuals, named Angelo, Bianca, Corin, Dolabella,
Emilia, and Fang, with salaries as follows:
Angelo
Dolabella
5000
9000
Bianca
Emilia
9000
5000
Corin
Fang
9000
11000
The average salary is $8000. However, $8000 is not the median salary, since only two
persons earn less, and four earn more. In this example, the median salary is $9000.
If three other persons are added to the sample, Grumio, Hamlet and Imogen, earning
$6000, $7000 and $8500, respectively, then the median salary becomes $8500.
Example 20.31 (Half-life and lifetime) Consider a radioactive particle that may disintegrate at
any moment. Start the experiment at time t = 0 (the particle has not yet disintegrated) and
assume that the probability that it does disintegrate during an infinitesimal time interval dt is
equal to dt (independently of t). As in Section 20.4.b, one can show that the probability that
the particle be intact at time t is given by
F (t) = e t .
The median of the random variable time before disintegration is called the half-life (the
probability of disintegration during this length of time is 12 ). To compute it, we must solve the
equation
1
F (t1/2 ) = exp{ t1/2 } = ,
2
and therefore we have t1/2 = (log 2)/. Note that the expectation of this r.v., which is called
the (average) lifetime of the particle, is equal to 1/. Indeed, the probability density for
disintegration at time t is given by
f (t) = F (t) = e t ,
and the expectation is therefore
Z
Z
t1/2
1
= t f (t) dt = t e t dt = =
> t1/2 .
log 2
It is such that F () = 1/e.
Random variables
534
20.6
Random vectors
20.6.a
Pair of random variables
Sometimes it is necessary to consider a random vector, or more generally an

n-tuple of random variables. We start with the case of a pair (X , Y ) of random
variables, all defined on the same probability space (, , P).
DEFINITION 20.32 The distribution function or joint distribution func-
tion of a pair (X , Y ) of random variables is the function

def
F X Y (x, y) = P{X x and Y y},

defined for x, y R, that is, the probability that both X is less than or equal
to x, and Y is less than or equal to y.
PROPOSITION 20.33 Let F X Y be the distribution function of the pair (X , Y ).
Then, for any x, y R, we have
F X Y (, y) = F X Y (x, ) = 0
and
F X Y (+, +) = 1.
Moreover, F X Y in increasing and right-continuous with respect to each variable.

The notion of probability density is also easy to extend:
DEFINITION 20.34 (Joint probability density) The probability density (also
called joint probability density) of the pair (X , Y ) is the function (or sometimes the distribution)
def
f X Y (x, y) =
We then have
F X Y (x, y) =
2 FX Y
(x, y).
x y
Z
f X Y (s, t) dt ds.
DEFINITION 20.35 (Marginal distribution) Let (X , Y ) be a pair of random

variables. The marginal distribution of X is the distribution of the random
variable X .
Knowing the (joint) distribution function of the pair (X , Y ), it is very easy

to find the distribution function of X , and hence its marginal distribution.
Since F X (x) is the probability that X x, with no condition whatsoever
imposed on the variable Y , we have
F X (x) = F (x, +),
Random vectors
535
and similarly the marginal distribution of Y is given by

F Y ( y) = F (+, y).
It follows that the marginal probability densities are
Z +
Z
f X (x) =
f X Y (x, y) dy
and
f Y ( y) =
f X Y (x, y) dx.
We now introduce the moments of a pair:

DEFINITION 20.36 Let (X , Y ) be a pair of r.v. and let k and be integers.
The ( k , )-th moment or moment of order ( k , ) of (X , T ) is the quantity

ZZ
def
k, =
x k y f X Y (x, y) dx dy,
R2
def
def
1,0 and m Y =
0,1 the expectations of X and
if it exists. If we denote m X =
Y , the corresponding centered moment is given by
ZZ
def
k, =
(x m X )k ( y m Y ) f X Y (x, y) dx dy.
R2
DEFINITION 20.37 The variances of (X , Y ) are the centered moments of
order (2, 0) and (0, 2), denoted Var(X ) = X2 = 20 and Var(Y ) = Y2 = 02 .

The covariance of (X , Y ), denoted Cov(X , Y ), is the centered moment of
order (1, 1), that is,
def
Cov(X , Y ) = 11 .
We then have the relation
h

i
Cov(X , Y ) = E X E(X ) Y E(Y ) .
Remark 20.38 Note also the obvious formula Cov(X , X ) = X2 = Var2 (X ).
Exercise 20.4 Let X and Y be two random variables with moments of order 1 and 2. Show
that
Var(X + Y ) = Var(X ) + Var(Y ) + 2 Cov(X , Y ).
Note that the centered moment of order (1, 1) is related with the noncentered moment by means of the formula 11 =
11 10 01 .
Finally, we have the correlation of random variables:
DEFINITION 20.39 The correlation of the pair (X , Y ) is the quantity
def
r=p
11
Cov(X , Y )
=p
20 02
Var(X ) Var(Y )
if the moments involved exist.
Random variables
536
This has the following properties (following from the Cauchy-Schwarz inequality):
THEOREM 20.40 Let (X , Y ) be a pair of random variables with correlation r,
assumed to exist. Then we have
i) |r| 1;
ii) |r| = 1 if and only if the random variables X and Y are related by a linear
transformation, more precisely if
X mX
Y mY
=r
Y
X
almost surely.
DEFINITION 20.41 (Uncorrelated r.v.) Two random variables X and Y are

uncorrelated if their correlation exists and is equal to zero.
The correlation is useful in particular when trying to quantify a link between two statistical series of numbers. Suppose that for N students, we know
their size Tn and their weight Pn , 1 n N , and we want to show that there
exists a statistical link between the size and the weight. Compute the average
weight P and the average size T . The standard deviations of those random
variables are given by
2 (T ) =
N
1 X
2
Tn2 T
N n=1
and
2 (P ) =
N
1 X
2
Pn2 P .
N n=1
The correlation between size and weight is then equal to

r=
N
X
1
(T T )(Pi P ).
(T ) (P ) n=1 i
If it is close to 1, there is a strong positive correlation between size and weight,

that is, the weight increases in general when the size increases. If the correlation is close to 1, the correlation is negative (the taller you are, the lighter).
Otherwise, there is little or no correlation. If |r| = 1, there is even an affine
formula relating the two statistical series.
The correlation must be treated with some care. A value r = 0.8 may
indicate an interesting link, or not, depending on the circumstances. For
instance, the size of the dataset may be crucial, as the following example
(taken from Jean Meeus [66]) indicates.
The table below gives the absolute magnitude5 m of some asteroids and the
number of letters in their name:
5
A logarithmic measure of brightness at fixed distance.
Random vectors
537
Name
Name
Niobe
Feronia
Klytia
Galatea
Eurydike
Freia
8.5
10.3
10.3
10.1
10.0
9.0
5
7
6
7
8
5
Frigga
Diana
Eurynome
Sappho
Terpsichore
Alkmene
9.6
9.1
9.3
9.3
9.7
9.4
6
5
8
6
11
7
The correlation for the first eight asteroids is r 0.785, which is quite large.
Worse, the correlation for the last five is r 0.932. Of course, this is merely
an accident, due to the small sample used. There is in fact no correlation
between the length of the name and the magnitude of asteroids. With the
twelve given, the correlation becomes 0.4 and with fifty asteroids6 it is only
r 0.097.
Remark 20.42 The covariance of (X , Y ) may be seen as a kind of scalar product of X E(X )
with Y E(Y ) and the standard deviation as the associated norm. The correlation is then
the cosine of the angle between the two functions in L2 space:
r=
(X |Y )
.
kX k2 kY k2
So there is no surprise in the fact that this number is between 1 and 1. Moreover, |r | = 1 if
and only if the functions X E(X ) and Y E(Y ) are linearly dependent, as we found above.
20.6.b Independent random variables

DEFINITION 20.43 (Independence) Two random variables X and Y are independent random variables if the joint distribution function can be expressed as the direct product (or tensor product) of their marginal distribution functions F X Y = F X F Y , that is, if
F X Y (x, y) = F X (x) F Y ( y) for all x, y R.
It is usual to abbreviate independent random variables by i.r.v.

There is an equivalent form of the definition which involves the -algebra
generated by a random variable:
DEFINITION 20.44 (-algebra generated by an r.v.) Let X be a random variable. The -algebra generated by X, denoted T (X ), is the set of all subsets
of the form X 1 (A) for a Borel subset A.
PROPOSITION 20.45 Two random variables X and Y are independent if and only
if the -algebras generated by X and Y are independent, that is, if, for all A T (X )
and B T (Y ), we have P(A B) = P(A) P(B).
6
Numbered 51 to 100 in the Ephemerids of Minor Planets (Leningrad, 1979).
Random variables
538
is,
If the joint probability density exists, we then have f X Y = f X f Y , that

f X Y (x, y) = f X (x) f Y ( y) for all x, y R.
We have the following important result:
THEOREM 20.46 Two independent r.v. X and Y are uncorrelated.

Proof. We can assume that the variances X2 and Y2 exist; otherwise the correlation
is not defined. Then we have
ZZ
Cov(X , Y ) =
(x m X )( y m Y ) f X Y (x, y) dx dy
Z
Z
= (x m X ) f X (x) dx ( y m Y ) f Y ( y) dy = 0,
since each factor vanishes.
The converse is not true in general, however.

Example 20.47 Let = [0, 2] with the uniform probability measure P = d/2.
Let
X () = cos and Y = sin . The r.v. X and Y are not independent (intuitively, because there
is a relation between them, namely X 2 + Y 2 = 1), but on the other hand the covariance is
equal to
Z
d
Cov(X , Y ) = cos sin
= 0,
2
so they are uncorrelated.
20.6.c Random vectors

It is now possible to generalize the preceding results to n random variables.
DEFINITION 20.48 An n-dimensional random vector is an n-tuple of ran-
dom variables, or in other words a measurable map

X = (X1 , . . . , Xn ) : Rn .
DEFINITION 20.49 The distribution function (or joint distribution func-
tion) of a random vector X is the function defined on Rn by

F X ( x) = P X i x i for i = 1, . . . , n ,
where x = (x1 , . . . , xn ).
PROPOSITION 20.50 Let F X be the distribution function of X. Then for all
x Rn we have
F X (, x2 , . . . , xn ) = = F X (x1 , . . . , , . . . , xn ) =
= F X (x1 , . . . , ) = 0
and
F X Y (+, . . . , +) = 1.
Morevoer, F X is increasing and right-continuous with respect to each variable.
Image measures
539
DEFINITION 20.51 The probability density (or joint probability density) of

a random vector X is the function (or distribution)
def
f X (x1 , . . . , xn ) =
We then have
F X (x1 , . . . , xn ) =
x1
n
F (x , . . . , xn ).
x1 . . . xn X 1
...
xn
f (s1 , . . . , sn ) ds1 dsn .
20.7
Image measures
20.7.a Case of a single random variable
Exercise 20.5 The probability density of the velocity of an atom of a perfect gas is known
(it is a maxwellian7). What is the probability density of the kinetic energy mv 2 /2 of this atom?
To treat this kind of problem, we must consider the situation where, given
a random variable X , a second random variable Y = (X ) is defined. This is
called a change of random variable. The task is to determine the distribution
function of Y , and its probability density, knowing those of X .
When the change of random variable is bijective, there is a unique y = (x)
associated to any x R, and it is then easy to see that if f denotes the
probability density of X and g denotes that of Y , then

f (x)
g( y) =
,
(x)
1
x=
( y)
provided has no zero.
Proof. Indeed, we have G ( y) = P{X 1 ( y)} = F 1 ( y) if is increasing, and

G ( y) = P{X 1 (x)} = 1 P{X < 1 ( y)} = 1 F 1 ( y) if is decreasing.
Then

f (x)
g( y) = G ( y) = F 1 = F 1 ( 1 ) = ,
(x)
as was to be shown.
If the function is not a bijection, we have the following more general

result (compare with the change of variable formula for a Dirac distribution,
Exercise 7.4.d, page 198):
7
In other words, a gaussian with a physicists name.
Random variables
540
THEOREM 20.52 Let : X 7 Y be a change of random variable and y R
such that the equation
y = (x)
has at most countably many solutions, and let S y denote the set of these solutions.
Assume, moreover, that (x) 6= 0 for x S y . Then the probability density g of Y is
related to the probability density f of X by
X f (x)

.
g( y) =
(x)
xS y
20.7.b Case of a random vector

Theorem 20.52 generalizes to a function : Rn Rn as follows:
THEOREM 20.53 Let : X 7 Y be a change of random vector, and let y Rn
be such that the equation y = ( x) has at most countably many solutions x k , k I .

Assume, moreover, that the differential d( xk ) is invertible for all i I . Then the
probability density g of Y is related to the probability density f ofX by

X
X
D(x1 , . . . , xn )
f ( xk )

,
=
f ( xk )
g( y) =
|det d(x )|
D( y , . . . , y )
k
kI
kI
where the last factor is the inverse of the jacobian matrix of the change of variable.
20.8
Expectation and characteristic function
20.8.a
Expectation of a function of random variables
Let X be a random variable and Y = (X ). The expectation of Y , if it exists,

can be computed using the probability density f X of X by the formula
Z
Z

E(Y ) = E (X ) =
X () dP() =
(x) f X (x) dx.
More generally, if Y = (X1 , . . . , Xn ) is a random variable defined by

change of variable from a random vector, the expectation of Y (if it exists)
may be computed using the joint probability density f X1 Xn (x1 , . . . , xn ) of
(X1 , . . . , Xn ) by the formula
ZZ
E(Y ) =
(x1 , . . . , xn ) f X1 Xn (x1 , . . . , xn ) dx1 . . . dxn .
Rn
Expectation and characteristic function
541
For instance, if X and Y are random variables, the expectation E(X Y ) of

the product X Y can be computed using the joint probability density f X Y :
Z
ZZ
E{X Y } =
X () Y () dP() =
x y f X Y (x, y) dx dy.
R2
THEOREM 20.54 The expectation has the following properties:
i) it is linear, that is, we have

nP
o P
E
i X i = i E{X i }
i
whenever the right-hand side is defined;
ii) ifX is almost

surely equal to a constant a, E{X } = a; in particular, we have

E E{X } = E{X };
iii) we have

E{X Y }2 E{X 2 } E{Y 2 },
by the Cauchy-Schwarz inequality, with equality if and only if X and Y are
linearly dependent;
iv) if X and Y are independent, then E{X Y } = E{X } E{Y }.
20.8.b Moments, variance

Using the formulae of the previous section, we see that the moments and centered moments of a random variable X are given by expectations of functions
of X , namely,
n
k o
k = E{X k }
and
k = E X E{X } .
Moreover, the k-th moment (centered or not) exists if and only if the function
x 7 x k f (x) is Lebesgue-integrable, where f is the probability density of X ,
that is, if
Z
k
x f (x) dx < +.
As a special case, the variance is given by

2 = E (X E{X })2 } = E{X 2 } E{X }2 .
20.8.c Characteristic function
DEFINITION 20.55 The characteristic function X of a random variable X

is the function defined for R by

X () = E e iX .
Random variables
542
The following result is then immediate.

THEOREM 20.56 The characteristic function X of a random variable X is the
Fourier transform8 of its probability density (in the sense of distributions):
Z
Z
iX
()
e
dP() =
e i x f (x) dx.
X () =
Notice that this quantity is always defined, for any random variable X , because
|e iX | = 1 and is of total measure equal to 1 (or equivalently, because the
density probability f , when it exists as a function, is Lebesgue-integrable).
According to the general theory of the Fourier transform, we have therefore:
PROPOSITION 20.57 Let X be a random variable and let X be its characteristic
function. Then

i) X (0) = 1 and X () 1 for all R;
ii) if the k-th order moment of X exists, X can be differentiated k times and we
have
(k)
k = E{X k } = (i)k X (0);
iii) since X is real-valued, X is hermitian, i.e., we have X () = X ().
R +
Proof. Note that X (0) = f X (x) dx = 1 or, to phrase it differently, we have
X (0) = E{1} = 1. The last two properties are direct translations of known results of
Fourier transforms.
THEOREM 20.58 The characteristic function of a random variable X determines
uniquely the distribution function of X .
In other words, two random variables have the same characteristic function if
and only if they have the same distribution function. (This theorem follows
from an inversion formula which is the analogue of the Fourier inversion
formula and is due to P. Lvy).
Before going to the next section, note also the following important fact:
PROPOSITION 20.59 The characteristic function of the sum of two independent
random variables is the product of the characteristic functions of the summands:
X +Y () = X () Y ().
8
Be careful with the various conventions for the Fourier transform! It is usual in probability
to omit the factor 2 which often occurs in analysis.
Sum and product of random variables
543
Proof. Let f X Y be the joint probability density of (X , Y ). Since X and Y are

independent, we have f X Y (x, y) = f X (x) fY ( y) and hence

X +Y () = E e i (X +Y )
ZZ
=
e i (x+ y) f X Y (x, y) dx dy
Z
Z
= e i x f X (x) dx e i y f Y ( y) dy = X () Y ()
by the expectation formula for a change of random variables.
20.8.d Generating function

DEFINITION 20.60 The generating function G X of a random variable X is
the function of a real variable s defined as the expectation of s X :
G X (s) = E(s X )
when it exists.
THEOREM 20.61 Let X : N be an integer-valued discrete r.v., and let pn =
P(X = n) for n N. Then G X is defined at least for s [1, 1] and is of C class

on ]1, 1[. Moreover, we have G X (1) = 1 and G X (s) 1 for all s [1, 1].
In addition, for s [1, 1], we have the series expansion
G X (s) =
In particular
E(X ) = G X (1)
n=0
(n)
pn
sn
and
and hence
G (0)
P(X = n) = X
.
n!

2
Var(X ) = G X (1) + G X (1) G X (1) .
20.9
20.9.a
Sum of random variables
Let us assume that every day I eat an amount of chocolate which is between
0 and 100 g, and let us model this situation by a random variable X1 for the
amount eaten on the first day, another random variable X2 for the amount
eaten on the second day, and so on, each random variable being assumed
to have a uniform distribution on [0, 100]. We would like to known the
distribution of the total amount of chocolate eaten in two days.
544
Random variables
Let us start with a slightly simplistic reasoning, where the amount of

chocolate can only take integral values k between 1 and 100 with probability P(X i = k) = 1/100. There is a very small probability that I will eat only
2 g of chocolate, since this requires that I eat 1 g the first day and 1 g the second, with a total probability equal to 1/10 000 (if we assume that the amount
eaten the first day is independent of the amount eaten the second). The probability that I eat 5 g is larger, since this may happen if I eat 1 g and then 4 g,
or 2 g and then 3 g, or 3 g and then 2 g, or (finally) 4 g and then 1 g; this gives
a total probability equal to 4/10 000.
The reader will now have no trouble showing that the probability that I
eat n grams of chocolate is given by the formula
P(X1 + X2 = n) = P(X1 = 1) P(X2 = n 1)
+ P(X1 = 2) P(X2 = n 2) + + P(X1 = n 1) P(X2 = 1).
With a finer subdivision of the interval [0, 100], a similar formula would hold,
and passing to the continuous limit, we expect to find that
Z 100
ftotal (x) =
f (t) f (x t) dt = f f (x),
0
where ftotal is the density of X1 + X2 and f that of X1 or X2 (since the support

of f is [0, 100]). In other words, the probability density of the sum of two
independent random variables is the convolution product of the probability
densities of the arguments.
This result can be established rigorously in two different manners, using a
change of variable or the characteristic function. The reader may find either
of those two proofs more enlightening than the other.
(Using the characteristic function)
From Proposition 20.59 we know that the characteristic function of the sum of
two independent random variables is the product of the characteristic functions
of the summands. Using the inverse Fourier transform, we obtain:
THEOREM 20.62 (Sum of r.v.) Let X1 and X2 be two independent random variables, with probability densities f1 and f2 , respectively (in the sense of distributions).
Then the probability density of the sum X1 + X2 is given by
f X1 +X2 = f1 f2 ,
or equivalently
f X1 +X2 (x) =
f1 (t) f2 (x t) dt,
which is only correct if the random variables involved are continuous, the probability
densities then being integrable functions.
545
If we denote by f2 the transpose of f2 , then the probability density of the difference

X1 X2 is given by
Z +
f X1 X2 = f1 f2
or
f X1 X2 (x) =
f1 (t) f2 (x + t) dt.
Note that
f X1 +X2 =
f X1
f X2 = 1, as it should be.
(With change of variables)
Let now f2 (x1 , x2 ) denote the joint probability density of the pair of random
variables (X1 , X2 ), and make the change of variables Y1 = X1 and Y2 =
X1 + X2 . The jacobian of this transformation is given by

D( y1 , y2 ) 1 0

= 1,
| J| =
=
D(x1 , x2 ) 1 1
and the joint probability density of (Y1 , Y2 ) is g( y1 , y2 ) = f2 ( y1 , y2 y1 ); the

probability density of Y = X1 + X2 = Y2 is the marginal density, so we get:
THEOREM 20.63 (Sum of r.v.) Let Y = X1 + X2 and let f2 (x1 , x2 ) denote the
joint probability density of the pair (X1 , X2 ). Then the probability density of Y is
given by
Z +
g( y) =
f2 (t, y t) dt.
If the random variables X1 and X2 are independent, we have f2 = f X1 f X2 and the

probability density of Y is the convolution g = f X1 f X2 :
Z +
g( y) =
f X1 (t) f X2 ( y t) dt.
Example 20.64 The sum of two independent gaussian random variables X1 and X2 , with means
m1 and m2 and variance 12 and 22 , respectively is another gaussian variable.

Adpating the result of Exercise 10.2 on page 295, which computes the convolution of two
centered gaussians, one can show that X1 + X2 is gaussian with mean m1 + m2 and variance
12 + 22 .
The following result is also very useful:

THEOREM 20.65 (Bienayme identity) Let X1 , . . . , Xn be independent random
variables, each of which has finite variance Var(X i ) = i2 . Then the variance of the
sum exists and is the sum of the variances of the X i : we have
Var(X1 + + Xn ) = Var(X1 ) + + Var(Xn ).
Proof. Without loss of generality, we may assume that each random variable is centered (otherwise, replace Xi by Xi E(Xi ) without changing its variance).
Random variables
546
Let i denote the characteristic function of the random variable Xi . We have

k (0) = 0 for k = 1, . . . , n, and by Theorem 20.59 on page 542 we can write
X1 ++Xn () = 1 () 2 () n ().
On the other hand, the variance of Xk is given by Var(Xk ) = k (0). It therefore suffices to differentiate the above function twice and then evaluate the resulting expression
at 0. Since the reader will easily check that all first derivatives give zero contributions,
there only remains the formula
X1 ++Xn (0) = 1 (0) 2 (0) n (0) + + 1 (0) n1 (0) n (0).
Since in addition we have k (0) = 1 for all k, the proof is complete.
20.9.b Product of random variables

While we are on this topic, we might as well consider the topic of product
of random variables. For this, put Y1 = X1 and Y2 = X1 X2 . This change of
variable has jacobian

D( y1 , y2 ) 1 0

= |x1 | ,
| J| =
=
D(x1 , x2 ) x2 x1
and hence we obtain:
PROPOSITION 20.66 (Product of r.v.) The probability density of Y = X1 X2 is
given by
y 1
dt,
f2 t,
t |t|
in terms of the joint probability density f2 of (X1 , X2 ). If the random variables X1

and X2 are independent, we have
Z +
y 1
g( y) =
dt.
f X1 (t) f X2
t |t|
g( y) =
Similarly, the reader will have no difficulty checking the following:

PROPOSITION 20.67 (Ratio of r.v.) The probability density of Y = X1 /X2 is
g( y) =
f2 (t, t y) |t| dt
and, if the random variables X1 and X2 are independent,

Z +
g( y) =
f X1 (t) f X2 (t y) |t| dt.
Bienaym-Tchebychev inequality
547
Irne Jules Bienaym (17961876), studied in Bruges and then at

the lyce Louis-le-Grand, and took part in the defense of Paris
in 1814. After going to the cole Polytechnique, his scholarly life
was perturbed by political turbulence (the school closed for some
time). He obtained a position at the military school of Saint-Cyr
and then became a general inspector in the Ministry for Finance.
After the 1848 revolution, he became professor of probability at
the Sorbonne and a member of the Acadmie des Sciences in 1852.
He had a knack for language and translated in French the works
of his friend Tchebychev. He also had a knack for disputing with
others, such as Cauchy or Poisson.
20.9.c Example: Poisson distribution

The following is easy to check:
THEOREM 20.68 Let X be an r.v. with a Poisson distribution with parameter .
Then the characteristic function of X is

X () = E e iX = exp (e i 1) .
Let X1 and X2 be two independent r.v. with Poisson distributions with parameters
1 and 2 , respectively. Then the r.v. X = X1 + X2 has a Poisson distribution with
parameter = 1 + 2 .
20.10
Bienaym
e-Tchebychev inequality
20.10.a Statement
Given a random variable X , its standard deviation mesaures the average
(quadratic) distance to the mean, that is, the probability that the actual value
of X be far from the expectation of X . The Bienaym-Tchebychev inequality
is a way to quantify this idea, giving an upper bound for the probability that
X differs from its expectation by a given amount.
THEOREM 20.69 (Bienayme-Tchebychev inequality) Let (, , P) be a probability space and let X be a real-valued random variable on such that E(X 2 ) < +.
Let m = E(X ) be the expectation of X and let 2 be its variance. Then, for any
> 0, we have

2

1
P |X m| 2 ,
or equivalently
P |X m| 2 .
Random variables
548
Pafnouti Tchebychev (18211894) entered Moscow University in

1837, and Saint Petersburg University in 1847. He studied prime
numbers (proving in particular that there always exists a prime
number p such that n < p 2n for all n 1, and giving the
first correct estimate for the order of magnitude of the number of
primes less a large quantity X , which was refined into the prime
number theorem in 1896 by Hadamard and de la Valle Poussin).
He was also interested in mechanics, in quadratic forms (Tchebychev polynomials), and in probability theory, where he published
four important papers. His works were described very clearly by
Andrei Markov. He died while drinking tea, like Euler, which is
rigorous proof that mathematicians should drink coffee.9
In the other direction, knowing the actual value of X for some experience
, this inequality gives an estimate of the error in assuming that this value is
equal to the average value m.
Example 20.70 Let X be a random variable with expectation m and standard deviation . The
probability that the result X () of an experiment differ from m by more than 4 is at most
6, 25 %.
To prove the inequality, we start with the following lemma:

LEMMA 20.71 (Markov inequality) Let (, , P) be a probability space and let X
be a non-negative random variable such that E(X ) < + (i.e., in the terminology of
analysis, we consider a measure space and a non-negative integrable function). For any
t R+ , we have

1Z
P {X t}
X () dP().
t
Proof. This is clear by looking at the following picture:
X ()
t
X () dP()
t P({X t})

The dark region is a rectangle with height t and width P {X t} . Of course, it
may well be something else than a rectangle (for instance many rectangles, depending
on the structure
of the event {X t}), but in any case the dark region
R has measure
t P {X t} . It is clear that this measure is smaller than or equal to X () dP(),
since X 0.
9
Mathematicians are machines for turning coffee into theorems, Paul Erds (or Alfred
Rnyi).
Bienaym-Tchebychev inequality
549
More rigorously, we can write

0 t {X t} () X ()
(since t {X t} () is equal to t or 0). Integrating over with respect to the probability
measure dP, the inequality follows.
2
Applying this lemma to X E(X ) / 2 and t = 2 , we derive the
Bienaym-Tchebychev inequality.
20.10.b Application: Buffons needle

Buffon10 computed in 1777 that when throwing a needle of length a onto a
floor with floorboards of width b a, the probability that the needle will
2a
cross a line between two boards is equal to b
. Take a = b to simplify, and
let p = 2/ denote the resulting probability.
Suppose the needle is thrown N times, where N is very large. If we
let X denote the random variable corresponding to the number of times
the needle has crossed a line between two planks, it follows that X has a
binomial distribution B(N , p). Thus the average number of needles crossing
a line is given by the expectation E(X ) = p N , and the standard deviation
of X is 2 = N p(1 p). We can expect than when N is very large, the
ratio X /N should be close to p = 2/, given an experimental numerical
approximation of (this experiment has been performed a number of times,
and many simulations can be found on the Web).
The question that may be asked is then: to obtain a approximation of
with a given precision, how many times should the needle be thrown? For
instance, suppose I want to obtain the value of with an error of at most
1/1000, how many times should I throw the needle to be certain that the
resulting value will be correct with a probability of 99%?
The simplest estimate for such a question can be obtained simply from the
Bienaym-Tchebychev inequality, which we recall for convenience:

1
P |X N p| 2 .
In order for the probability that X satisfies the inequality

to be at most 1%,
p
we must take 1/2 = 1/100; hence we take = 100 = 10. Let be the
value of obtained experimentally, namely = 2N /X . We want to
satisfy 103 + 103 , and therefore for the number X of needles
crossing a line to satisfy
2N
2N
.
+ 103
103
10
Georges Leclerc, comte de Buffon (17071788), famous naturalist, author of the monumental Histoire naturelle, was also interested in the theory of probability.
Random variables
550
Writing |X N p| = with values

X =
2N
103
(both signs + and lead to the same result), we find

N=
2 p
N p(1 p),
2 103
hence
N=
4 2
p(1 p)
4 106
in other words, N = 560 million, more or less. It is not very reasonable to

expect that can be computed in this manner.
However, this is a very coarse estimate, since the Bienaym-Tchebychev
inequality holds with no assumptions on the distribution of the random
variable being considered. If we take into account the precise distribution of
X , a more precise result can be derived. But this remains of the same order of
magnitude, and does not really affect the discussion. For instance, since the
binomial distribution may be approximated by a gaussian distribution for N
such that N p(1 p) is very large, we can find numerically (using tables or
computations of the function erf(x)) a value N 150 million throws, which
is only better by a factor of 3.
20.11
Independance, correlation, causality
It is very important to distinguish between independance, correlation, and
causality. In Example 20.47, we saw that two dependent events can be uncorrelated.
Two events are causally linked if one is the cause of the other. For instance,
statistics show that road accidents are more frequent on Saturdays than other
days of the week (thus they are correlated). There is causality here: usually
more alcohol is consumed on Saturday, leading to a higher rate of accidents.
However, events may be correlated without causality. An example is given
by Henri Broch [15]: in villages in Alsace, statistics show that the number of
births per year is highly correlated with the number of storks (the correlation
is close to 1). Should we conclude that storks carry babies home (which would
mean causality)?11
David Ruelle, in his excellent book [77], mentions events which are causally
linked but uncorrelated, which is more surprising. Consider, on the one hand,
the position of the planet Venus in the sky, and on the other hand the weather
11
There is in fact a hidden causal explanation: the more families in a village, the more
houses, hence the more chimneys, and the more room for storks. This causality is not because
of the two events discussed, but because of a third, linked independently with both.
Independance, correlation, causality
551
one month after the date of observation. We know very well that weather [86]
is a highly chaotic system, and hence is very sensitive to initial conditions. The
position of Venus on a given day is therefore highly likely to influence the
weather one month later. There is therefore strong causality (and dependence)
between the position of Venus and the weather. However, no correlation can
be observed. Indeed, other factors (the day of the year, the weather in previous
months, the solar activity, the position of other planets, etc.) are equally or
more crucial in the evolution of the weather than the position of Venus.
Looking only at Venus and averaging over all other variables, no trace of the
influence of this fair planet will be noticeable. Causality without correlation.
552
Random variables
Sidebar 8 (Riemann-Stieltjes integral) Let : R R be an increasing func-
tion (not necessarily

continuous). A measure is defined on open intervals by

]a, b[ = (b ) (a + ) (where (b ) denotes the limit of (x) as x tends
to b on the left, and (a + ), similarly, is the limit as x tends to a on the right; both
are well defined because is increasing). This measure may be extended to Borel sets
in R uniquely, by defining
S P
i Bi =
i (B i )
if B i B j = for i 6= j. This yields a Borel measure, that is, a measure on

(X , B), with the property that

[a, b] = (b + ) (a ),
[a, b[ = (b ) (a ), etc.
The integral of a functionR f with respect

R to this measure is defined as described in
this chapter, and is denoted f d or f d, and called the Riemann-Stieltjes
integral of f with respect to . It satisfies the properties
stated in Theorem

2.28 on
page 65. There are associated spaces L1 R, d and L1 [a, b] , d , which are
complete by the Riesz-Fischer theorem.
Finally, if is of C 1 class, one can show that
Z
Z
Z
d
f d = f (x) (x) dx = f (x)
dx
dx
for any function f integrable with respect to .

This construction is interesting in part because it gives a definition of certain
integrals or generalized functions without requiring the theory of distributions.
Example 20.72 Let (x) = 1 for x 0 and (x) = 0 otherwise. This is an increasing
function. The associated measure is called the Dirac

measure; it has the property
R
that any Dirac-measurable function is integrable and f d = f (0) for any f . Note that
functions in L1 (R, d) are always defined up to equality almost everywhere, but there {0}
is of measure 1 for , so it is not a negligible set.

Let P = x R ; {x} 6= 0 (the set of discontinuities of ). One can
define a measure p by p (X ) = (P X ), and moreover c = p is
also a Borel measure.
Chapter
21
Convergence of
random variables:
central limit theorem
Guildenstern (Flips a coin): The law of averages, if I have got this
right, means that if six monkeys were thrown up in the air for
long enough they would land on their tails about as often as they
would land on their
Rosencrantz: Heads. (He picks up the coin.)
Tom Stoppard, Rosencrantz & Guildenstern are dead [87]
21.1
Various types of convergence
We are going to define three types of convergence of a sequence (Xn )nN
of random variables to a limit X .
The first two are parallel to well-known concepts in the theory of integration.
Let (X , T , ) be a measure space (for instance, X = R, T = B, the Borel
-algebra, and the Lebesgue measure). A sequence ( fn )nN of measurable
functions on X converges almost everywhere to a measurable function f
if there exists a measurable subset
N X such that (N ) = 0 and, for
all x
/ N , the sequence fn (x) nN converges to f (x). In other words,
Central limit theorem
554
the sequence ( fn )nN converges simply to f for almost all x. In probability

language, we have the following analogous definition:
DEFINITION 21.1 (Almost sure convergence) Let (, , P) be a probability
space, (Xn )nN a sequence of random variables, and X a random variable.

c.a.s.
Then (X n ) nN
N converges almost surely to X, denoted X n X, if there
exists a subset N with probability zero such that, for any
/ N , the
sequence Xn () nN converges to X .
In integration theory, another notion is convergence in measure: a sequence
( fn )nN of measurable functions converges in measure to a measurable function f if
n
o

> 0
lim x X ; f (x) fn (x)
= 0.
n
In other words, with a margin of error > 0, the subset of those points where
fn is too far from f shrinks down to nothingness as n becomes larger and
larger.
Similarly, we have the probabilistic analogue:
DEFINITION 21.2 (Convergence in probability) Let (Xn )nN be a sequence of
r.v. and let X be a r.v. The sequence (X n ) nN

N converges in probability to X,
c.p.
denoted X n X, if
n
o

> 0
lim P ; X () Xn ()
= 0.
n
In general, almost everywhere convergence does not imply convergence in

measure. However, when the measure space has finite total measure ((X ) <
+), this is the case (for instance, this follows from Egorovs theorem, stated
on page 69). The case of probability spaces is of this type, since P() = 1, so
we deduce:
THEOREM 21.3 Almost sure convergence implies convergence in probability.
The third type of convergence, which is very important in probability, is

the following:
DEFINITION 21.4 (Convergence in distribution) A sequence (Xn )nN of r.v.
with distribution functions (Fn )nN converges in distribution if there exists

a distribution function F such that the sequence (Fn )nN converges simply
to F at any point of continuity of F . This function determines the limit
distribution of the sequence (Xn )nN .
If X is an r.v. with distribution function F , then we say that (X n ) nN
N
c.d.
converges in distribution to X, and we denote X n X.
The law of large numbers
555
c.d.
Note that if Xn F , then we have
P(Xn x) P(X x)
n
for any x R which is a point of continuity of F .

Convergence in distribution is the weakest of the three types of convergence introduced. In particular, we have the implications:
c.a.s.
c.p.
c.d.
(Xn X ) = (Xn X ) = (Xn X ).

Certainly, convergence in distribution does not imply convergence in probability or almost sure, since the random variables Xn and X may have nothing to
do with each other; in fact, they may even be defined on different probability
spaces, since only the distribution function matters! It is perfectly possible to
have two r.v. X and Y with the same distribution, but which are never equal
(see page 524).
21.2
The law of large numbers
Consider a game of dce. I roll the die, and you must announce the result
in advance (with money at stake). If I only roll the die once, there is nothing
interesting you can really say (except if my die was false, but I would not say
so in advance). However, if I roll the die ten times, and you must guess the
sum of the results of the ten throws, there is much more you can say. More
precisely, you may resonably think that the sum will be close to 35, maybe
33 or 38, but most certainly not 6 (there is only one chance in 610 , that is,
one in sixty million to obtain a 1 ten times in a row; you may as well play
the lottery). And as I roll the die more and more times, you will be able to
bet with increasing confidence that the sum will be roughly n 3.5. This, in
effect, is what the weak law of large numbers makes precise1:
Let (, , P) be a probability
space and let (Xn )nN be a sequence of square integrable independent r.v. with the
same expectation m (i.e., for all n, we have E(Xn ) = m) and the same variance 2 .
Define
X + + Xn
Yn = 1
.
n
THEOREM 21.5 (Weak law of large numbers)
This theorem was proved (rigorously) for the first time by the Swiss mathematician James
(Jakob) Bernoulli (16541705) (who also discovered the famous Bernoulli numbers), in his
treatise Ars Conjectandi, published in 1713. His nephew Daniel Bernoulli (17001782) also
studied probability theory, as well as hydrodynamics. A generalization of Bernoullis result was
given by the marquis de Laplace, in his Thorie analytique des probabilits.
556
Then the sequence (Yn )nN converges in probability to the constant random variable
equal to m (i.e, Y () = m for all ). In other words, we have

> 0
P |Yn m| 0.
n
Proof. The simplest proof uses the Bienaym-Tchebychev inequality: notice that
E(Yn ) = m and then write
2 (Yn )
P |Yn m|
.
2
Since the random variables X1 , . . . , Xn are independent, the Bienaym identity 20.65
yields the formula
2 (n Yn ) = 2 (X1 + + Xn ) = n 2 ,
2 (Yn ) =
hence
from which we derive
n 2
2
=
,
2
n
n

2
P |Yn m| 2 0.
n n
Remark 21.6 It is possible to show that this result is still valid without assuming that the vari-
ables are square integrable (but they are always integrable). The proof is much more involved,
however.
There is a stronger version of the law of large numbers, due to Kolmogorov

and Khintchine (see, e.g., [33]):
THEOREM 21.7 (Strong law of large numbers) Let (Xn )nN be a sequence of in-
dependent identically distributed r.v. Assume that the Xn are integrable, and let m
denote their common expectation. Then the sequence of Cesro means converges almost
surely to m: we have
X1 + + Xn c.a.s.
m.
n
21.3
Let us come back to our game of dice. We know, according to the law of
large numbers, that we must bet that the sum is close to the average. It would
be interesting to know how the sequence (Yn )nN of Cesro means converges
to the constant random variable m. And more precisely, what should we expect
if we bet on a value other than the average?
2
He showed in particular that it is possible to cut a cake into 65, 537 equal pieces with
straightedge, compass, and knife. A very useful trick if you plan a large wedding.
557
A child prodigy, Carl Friedrich Gauss (17771855), son of an artisan,

entered the University of Gttingen at eighteen. He made contributions to astronomy (rediscovering Ceres, which lesser minds had
managed to lose in the firmament, by means of orbit computations)
and geodesy. With Wilhelm Weber, he created the theory of magnetism. However, his greatest discoveries were in mathematics. At
nineteen, he proved that the regular seventeen-sided polygon can
be constructed with straightedge and compass, and further generalized this result.2 One of the greatest among arithmeticians (he
introduced congruences, for instance), he also laid the foundation
of differential geometry, defining in particular the Gauss curvature
of a surface. He was named the prince of mathematicians.
The fundamental result (in particular, in many physical applications) is

that the successive sums of the terms of a sequence of independent, identically distributed random variables (square integrable), will always converge (in
distribution) to a Gaussian random variable, after an appropriate scaling. In
particular, in statistical physics, this is why many measurable quantities can be
expressed by the computation of a moment of Gaussian variables. Problem 8
gives the example of a simple problem of a random walk in the plane.
The theorem stating this fact is called the central limit theorem.3 In the
case of Bernoulli random variables (heads or tails), it was stated and proved
by Abraham de Moivre.
We begin by recalling the main facts conerning the Gaussian or normal
distribution, which is so essential in this result.
DEFINITION 21.8 (Normal distribution) The normal distribution with mean
m and standard deviation > 0, or Gaussian distribution with mean m and
standard deviation > 0, denoted N (m, 2 ), is the probability distribution
on R with density given by

1
(x m)2
gm, (x) = p
exp
.
2 2
2
A normal distribution of the type N (0, 2 ) is called a centered normal

distribution.
The reader will easily check that the characteristic function of the normal
distribution N (m, 2 ) is given by
2 2
2 2
u
u
m, (u) = e i um exp
, in particular, 0, (u) = exp
.
2
2
3
Translated lierally from the German, where a more suitable translation would have been
fundamental limit theorem (of probability theory).
558
Paul Lvy (18861971), son and grandson of mathematicians, studied at the cole Polytechnique, then at the cole des Mines, where
he became professor in 1913 . Hadamard, one of his teachers,
asked him to collect the works of Ren Gteaux, a young mathematician killed at the beginning of the First World War. After
doing this, Lvy developed the ideas of Gteaux on differential
calculus and functional analysis. At the same time as Khintchine, he studied random variables, and in particular the problem of the convergence of sequences of random variables. He
also improved and expanded considerably the results of Wiener
concerning Brownian motion. He was also interested in partial
differential equations and the use of the Laplace transform.
THEOREM 21.9 (Central limit theorem) Let (Xn )nN be a sequence of indepen-
dent, identically distributed random variables with finite expectation m and finite
variance 2 . Let
p
def X + + Xn
Yn = 1
and
Zn = n(Yn m).
n
Then the sequence (Zn )nN converges in distribution to the centered normal distribution N (0, 2 ). In other words, we have

Zx
Sn n
2
P
e t /2 dt.
p < x N (x) =
n
n
for all x R.
def
Proof. Let Sn = X1 + + Xn . We can compute the characteristic function of Zn ,

namely:
Sn p

i u p i u nm
Zn (u) = E e i u Zn = E e n

n

p
p
u
u
= e i um n X p
,
= e i um n Sn p
n
n
where the last relation follows from the independence of the r.v. Xn and from Theorem 20.59.
p
We now perform a Taylor expansion of X u/ n to order 2, using the fact that
and this yields
X (0) = iE(X )
and
X (0) = E(X 2 ),

2
u
u
u2
u
X p
=1+ip m
E(X 2 ) + o
.
n
n
2n
n
Inserted in the preceding relation, this implies that

u22
1
Zn (u) = exp i
+o
,
2
n
so the characteristic functions of Zn converge pointwise to that of the normal distribution. This turns out to be enough to imply the result (which is not too surprising, since
the characteristic function of an r.v. determines its distribution uniquely, as recalled
earlier), but the proof of this implication involves the continuity theorem of Paul Lvy
and is by no means easy.
559
Remark 21.10 The condition that the variance Var(Xn ) be finite is necessary. Exercise 21.15
on page 562 gives a counterexample to the central limit theorem when this assumption is
not satisfied. It should be mentioned that non-Gaussian processes are also very much in current
fashion in probability theory.
Application
Consider a piece of copper, of macroscopic size (electrical wiring, for instance).
Because of thermal fluctuations,
the average velocity of a conducting electron
p
in this copper wire is 3kb T /m, that is, of the order of 105 m s1 , or tens
of thousands of kilometers per second! Yet, the copper wire, on a table, does not
move. This is due to the fact that the velocities of the electrons (and of nuclei,
although they are much slower) are randomly distributed and cancel each
other. How does this work?
Let N be the number of conducting electrons in the wire. We assume that
the wire also contains exactly N atoms,4 and we will denote by m and M ,
respectively, the masses of the electrons and atoms. You should first check
that the velocity of atoms is negligible in this problem.
We denote by Xn the random variable giving the velocity (at a given time)
of the n-th electron. What is called the average velocity is not the expectation
of Xn , since (in the absence of an electric field), the system is isotropic and
E(Xn ) = 0. In fact, the average velocity is best measured by the standard
deviation of the velocity (Xn ).
What is then average velocity (the standard deviation of the velocity) of
the whole wire? Since the number N of atoms and electrons is very large (of
the order of 1023 ), we can write
N
n
X
m
1 m
1 X
Y =
X p
p
Xn ,
N (m + M ) n=1 n
N M
N n=1
and the term in parentheses is close to a centered gaussian with variance 2
by the central limit theorem.
The standard deviation of Y is therefore very close to
1 m
(Y ) = p
.
N M
p
Notice the occurrence of the famous ratio 1/ N , which is characteristic of
the fluctuations of a system with N particles. In the present case, this reduces
the average velocity of the set of electrons to 1.5 107 m s1 .
4
Depending on the metal used, it may not be N , but rather N /2 or N /3; however, the
principle of the calculation remains the same.
560

EXERCISES
Combinatorics
Exercise 21.1 All 52 cards in a pack of cards are dealt to four players. What is the probability
that each player gets an ace? Does the probability change a lot when playing algebraic whist?
(Recall that algebraic whist is played with a pack of 57 cards.)
Hint: Compute first, in general, the number of ways of subdividing n objects into k
subsets containing r1 , . . . , rk objects, respectively.
Exercise 21.2 Take a pack of n cards, and deal them face up on a table, noting the order the
cards come in. Then shuffle the pack thoroughly and deal them again. What is the probability
that at least one card occurs at the same position both times? Does this probability have a
limit as n goes to infinity?
It is amusing to know that this result was used by con-men on Mississippi steamboats, who
exploited the passengers intuitive (wrong) assumption that the probability would be much
smaller than it actually is.
Conditional probabilities
Exercise 21.3 Mrs. Fitzmaurice and Mrs. Fitzsimmons discuss their respective children
during a plane trip. Mrs. Fitzmaurice says: I have two children, let me show you their
pictures. She looks in her bag, but finds only one picture: This is my little Alicia, playing
the piano last summer.
What is the probability that Alicia has a brother?
Exercise 21.4 A religious sect practices population control by means of the following rules:
each family has children (made in the dark) until a boy is born. Then, and only then, they
abstain for the rest of their days. If the (small) possibility of twins is neglected, and knowing
that without any control, about 51% of babies are boys, what will be the percentage of males
in the community?
Same question if each family also limits to twelve the total number of children it can have.
Exercise 21.5 A new diagnostic test for a lethal disease has been devised. Its reliability is
impressive, says the company producing it: clinical trials show that the test is positive for
95% of subjects who have the disease, and is negative for 95% of subjects who do not have it.
Knowing that this disease affects roughly 1% of the population, what probability is there that
a citizen taking the test and getting a positive result will soon end up in eternal rest in the
family plot?
Exercise 21.6 Another amusing exercise. The other day, my good friend J.-M. told me that
he knows someone to whom the most amazing thing happened. This other person was just
thinking of one of his former relations, who he had not thought about for 10 years, and the
very next day he learned the person had just died, at the same time he was thinking of him!
(Or maybe within ten minutes.)
What should one think of this extraordinary (rigorously true) fact? Does this kind of
thing happen often, or do we have here an example of ESP (Extra-Sensory Perception, for the
uninitiated)?
Exercise 21.7 In 1761, Thomas Bayes, Protestant theologian, leaves for ever this vale of tears.
He arrives at the Gates of Heaven. But there is not much room left, and anyway God is a strict
Presbyterian and Bayes, unluckily for him, was a Nonconformist minister. However, St. Peter
Exercises
561
gives him a chance: Bayes is shown three identical doors, two of which lead to Hell and only
one to Heaven. He has to choose one. With no information available, Bayes selects one at
random. But before he can open it, St. Peter who is good and likes mathematicians stops
him. Wait, here is a clue, he says, and opens one of the two remaining doors, which Bayes
can see leads to Hell.
What should Bayes do? Stick to his door, or change his mind and select the other unopened
door?5
Expectation, standard deviation

Exercise 21.8 When a students hands in an exam paper, it will get a grade corresponding
roughly to the objective quality level of the paper. However, because of random perturbations
of the graders environment (did he eat well at lunch, is his cat purring on his knees,...), the
actual grade is really a random variable, with expectation equal to the objective grade m, and
with standard deviation = 2, say, for a paper graded out of 20 points.
To get a more exact result, it is possible to ask N different graders to look at the paper
and then average their respective grades, each of which is a random variable with the same
distribution, with expectation m and standard deviation .
Denote by X1 , . . . , X N the corresponding random variables for each grader, and let S N =
(X1 + + X N )/N be the average of their grades. What is the expectation of S N ? What is its
standard deviation?
How many correctors are needed so that the probability is at least 0.8 that the average grade
corresponds to the exact grade, within half a point? Same question if the standard deviation
is only half a point. Do you have any comments about the results?
Exercise 21.9 Solveig wants to make a chocolate cake. The recipe states that two large eggs
(100 g) are needed. But she only has two eggs from two different brands. The first package
says the eggs inside are large (between 50 and 60 g, with mass uniformly distributed), and the
second says the eggs are medium (between 40 and 45 g with uniform distribution). What is
the probability density of the total mass of the two eggs? What is the average and standard
deviation of this mass? What is the probability that the mass is greater than 100 g?
Covariance
Exercise 21.10 Let (Xn )nN be a sequence of random variables, pairwise independent, with
the Bernoulli distribution.6 with parameter p. Let Yn = Xn Xn+1 for n N .

1. Find the distribution of Yn .

2. Compute Cov Yn , Yn+k for n, k N .
3. Let
Sn =
n
X
Yi
n
i=1
for n N . Show that the sequence (Sn )nN converges in probability to the determinist
(constant) random variable p 2 .
This version of the classical Monty Hall paradox was suggested by Claude Garcia.
Recall that a Bernoulli distribution with parameter p is the distribution of a random
variable which takes the value 1 with probability p and 0 with probability 1 p.
6
562
Sum and product of r.v.

Exercise 21.11 Show that the ratio of two independent centered normal random variables
is a r.v. with a Cauchy distribution; recall that this means that for some a R, the probability
density is
a
1
f (x) =
.
a2 + x 2
Exercise 21.12 Let X be Y be two independent random variables, such that X has a Poisson
distribution with parameter and Y has a Poisson distribution with parameter , where
, > 0. Let pn = P(X = n) and qn = P(Y = n) for n N.
i) Show that Z = X + Y has a Poisson distribution with parameter + .
ii) Let k N. Compute the conditional probability distribution of X knowing that

Z = k.
iii) Let W = X Y . Compute the expectation and the variance of W .
Convergence
Exercise 21.13 A player flips two coins B1 and B2 repeatedly according to the following
rule: if flipping whichever coin is in his hand yields tails, switch to the other coin for the next
toss, otherwise continue with the same. Assume that tails is the outcome of flipping B1 with
probability p1 = 0.35, and of flipping B2 with probability p2 = 0.4. Let Xk be the random
variable equal to 1 if the player is flipping coin B1 at the k-th step, and equal to 0 otherwise.
Let
h (k)
h1 (k) = P(Xk = 1),
h2 (k) = P(Xk = 0),
and
H (k) = 1
.
h2 (k)
i) Show that H (k + 1) = M H (k), where M is a certain 2 2 matrix. Deduce that
H (k) = M k H (0), where H (0) is the vector corresponding to the start of the game.
ii) Diagonalize the matrix M . Show that as k tends to infinity, H (k) converges to a vector
which is independent of the initial conditions. Deduce that the sequence (Xk )kN
converges almost surely to a random variable, and write down the distribution of this
random variable.
Exercise 21.14 Let (Un ) be a sequence of identically distributed independent random vari-
ables such that
P(Un = 1) = q = 1 p with p [0, 1] .

Q
Define a sequence of random variables Vn by Vn = ni=1 Ui . Find the distribution of Vn , then
show that (Vn ) converges in distribution to a random variable V , and find the corresponding
distribution (it may be helpful to compute the expectation of Vn ).
P(Un = 1) = p,
Exercise 21.15 (Counterexample to the central limit theorem) Show that if X1 and X2 are
Cauchy random variables with parameters a and b , respectively (see Exercise 21.11), then the
random variable X1 + X2 is a Cauchy random variable with parameter a + b (compute the
characteristic function of a Cauchy random variable). What conclusion can you draw regarding
the central limit theorem?
Exercises
563
PROBLEMS
Problem 7 (Chain reaction) Let X and N be two square integrable random variables taking
values in N. Let (Xk )k1 be a sequence of independent random variables with the same distribution as X . Assume further that N is independent of the sequence (Xk )k1 . Let GX and G N
denote the generating functions of X and N . For n 1, let
Sn =
n
X
Xk .
k=1
Finally, let S be defined by

S () =
N
()
X
Xk ().
k=1
i) For n N, compute the generating function Gn of Sn , in terms of GX .

ii) Show that GS (t) = G N GX (t) = G N GX (t) for any t [1, 1].
iii) Compute E(S ) and V (S ) in terms of the expectation and the variance of X and N .
Consider now a reaction (for instance, a nuclear reaction) occurring in a certain

environment. Time is measured by an integer n 0 for simplicity. A particle (time 0)
comes into the environment, is absorbed, and produces new particles (time 1). In turn,
each new particle generates new particles at the next instant.
We denote by p0 the (constant) probability that a particle creates no other particle,
and for k N , we let pk be the probability that it creates exactly k particles. Let
t R
G (t) =
pk t k .
k=0
Finally, let Zk be the random variable which is the total number of particles at time k,
and let xn be the probability that no particle survives at time n (end of the reaction):
n N
xn = P(Zn = 0).
iv) Find G when p0 = 1. We now assume that p0 < 1. Show that G is strictly increasing
on [0, 1].
v) Find G if p0 + p1 = 1. Assuming now that p0 + p1 < 1, show that G is convex.
vi) In the following cases, find the number of solutions to the equation G (x) = x:
a) p0 = 0 ;
b ) 0 < p0 < 1 and G (1) 1 ;
c ) 0 < p0 < 1 and G (1) > 1
vii) Find a recurrence relation expressing the generating function Gn+1 of Zn+1 in terms of
the generating function Gn of Zn and of G .
viii) Still under the condition 0 < p0 < 1, find a recurrence relation between xn and xn+1 .
Show that the sequence (xk )kN is increasing, and deduce that it converges to a limit
such that 0 < < 1.
ix) Show that E(Zn ) = E(Z1 )n .
x) Show from the preceding results that:
(a) if E(Z1 ) 1, then the probability that the reaction stop before time n converges
to 1 as n tends to infinity;
(b) if E(Z1 ) > 1, then this probability converges to a limite 0 < < 1. What should
we expect from the reaction? What are your comments?
564
Problem 8 (Random walk in the plane) Consider a plane named P , a point named O which
serves as origin for some cartesian coordinate system on P , and a drunk mathematician named
Pierre who, having profited extensively from his pubs extensive selection of single malt whiskey,
leaves the afore mentioned pub (located at the origin O ) at time t = 0 to walk home. Each
second, he make a single step of fixed unit length, in a random direction (obstacles, constables,
and other mishaps are considered negligible).
Let Xn and Yn be the r.v. corresponding to the x- and y-coordinates of Pierre after the
n-th step, and let n denote the angle (measured from the x-axis) giving the direction following
during the n-th step. Thus we have
n
n
X
X
Xn =
cos k
and
Yn =
sin k .
k=1
k=1
We assume that the k are independent random variables uniformly distributed between 0 and
2.
The objective is to study the r.v.
Rn = Xn2 + Yn2
representing the distance from the pub.
Show that the r.v. Xn and Yn have expectation zero, and compute their variance.
Are X and Y correlated? Are they independent?
What is the expectation of Rn2 ?
Explain why it is reasonable to assume that for large n, the joint distribution of (Xn , Yn )
may be approximated by the joint distribution of a pair of independent centered normal
distributions. (Use the fact that uncorrelated gaussian variables are in fact independent.)
What is the (approximate) probability density of the pair (Xn , Yn )?
v) Deduce the (approximate) distribution function, and then the probability density,
of Rn .
vi) Finally, find an approximation for the expectation E(Rn ).
i)
ii)
iii)
iv)
SOLUTIONS
Solution of exercise 21.3. The probability that Alicia has a brother is 23 , and not 12 .
Indeed, among families with two children, about 25% have two boys, 25% have two girls,
and therefore 50% have one each. Since we know that there is at least one girl, the probability
of having two girls is 41 divided by 12 + 14 , which is 13 .
If Mrs. Fitzmaurice had distinguished between her children (saying, for instance, Here
is my older child at the piano), the probability would have been 12 . (Among those families
whose first child is a girl, there are 50% where the second is a girl.)
Solution of exercise 21.2. Let Ai denote the event which is the i-th card is at the same
place both times. It is easy to see that P(Ai ) = 1/n, that P(Ai A j ) = 1/n(n 1) for i 6= j
and that more generally
(n k)!
P(Ai1 Aik ) =
n!
when all indices are distinct. Using the Poincar formula, the probability of the event B: at
least one of the cards is at the same place both times is
!
n
[
P(B) = P
Ai ,
i=1
565
and is therefore equal to

P(B) =
There are
n
k
n
X
(1)k+1
k=1
1i1 <i2 <<ik
(n k)!
.
n!
n
distinct k-tuples in the sum, and we obtain

P(B) =
n
X
(1)k+1 Cnk
k=1
n
X
(n k)!
(1)k
= 1
.
n!
k!
k=0
This probability converges to 1 1/e (this is the beginning of the power series expansion for
e 1 ), or roughly 63%.
Solution of exercise 21.4. The probability hasnt changed: there are still 51% of baby boys
in the community.
This is easy to check: notice that every family has exactly one boy; those with two children
have exactly one boy and one girl, those with three have one boy and two girls, etc. The families
with n children have one boy and n 1 girls.
Let P (n) be the probability that a family has n children. We have of course P (n) = p q n1 ,
with p = 0.51 the probability that a baby is a boy, and with q = 1 p.
The proportion of boys is then given by
t=
P (1) + P (2) + + P (n) +

p(1 + q + q 2 + )
p/(1 q)
=
=
= p.
P (1) + 2P (2) + + n P (n) +
p(1 + 2q + 3q 2 + )
p/(1 q)2
Solution of exercise 21.5. Let P denote the event the result of the test is positive, and let
P be the opposite event. Let M denote the event the person taking the test has the disease,
and let M be its opposite.
From the statement of the problem, we know that P(P |M ) = 0.95 and that P(P |M ) = 0.95.
Since (M , M ) is a partition of the sample space (either a person has the disease, or not), by
Bayess formula, we have
P(M |P ) =
P(P |M ) P(M )
P(P |M ) P(M ) + P(P |M ) P(M )
With P(P |M ) = 1 P(P |M ) = 0.05, the probability of having the disease when the test is
positive is only around 16%.
Solution of exercise 21.6. Before making computations, two comments are in order. Saying
that the story is true does not mean that premonition really occurred, but that this is how the
story was told. Moreover, you are entitled to disbelief it is not because it is printed, with True
Story explicitly stated, even in Bold Face, that it is really true. For this type of debate, this is
an important remark.
What is the probability for one given person, whom I think about once in ten years, that I
will do so within ten minutes of his death? Lets spread out death on an interval of 80 years,
or roughly 42 million minutes, so that we can guess that the probability is roughly one out of
4.2 million.
I estimate that I know, more or less, roughly a thousand persons (including those I only
think about only once ten or twenty years); there may therefore be about one million people
who know someone I know.7 Hence it is not only possible, but in fact rather unsurprising,
that someone could tell me such a story!
7
This may be of the wrong order of magnitude, since my acquaintances each have, among
the people they know (those they only think about once in ten years, including people
known only by name, actors, and so on) rather more than a thousand individuals. On the
other hand, the sets are certainly not disjoint. Any finer estimate would be interesting.
566
This does not mean that ESP does not exist; only that there is no need here to invoke a
paranormal phenomenon in order to explain what happened.
Solution of exercise 21.7. Let us assume that door A leads to Heaven, and doors B and
C to Hell. If Bayes had chosen A, Saint Peter will open B or C for him; changing his mind
means Bayes will go to Hell.
If Bayes had chosen door B, Saint Peter, necessarily, must open door C for him (since he
opens a door leading to Hell), so if Bayes changes his mind, he will go to Heaven. The same is
obviously true if Bayes had chosen door C .
Therefore, independently of his original choice, if Bayes changes his mind after Saint Peter
has shown him one door, this change of mind will switch his destination from the one first
chosen. This means the probability that Bayes will go to Heaven is now 2/3 and the probability
of going to Hell is now 1/3.
Conclusion: well, this rather depends on what Bayes really wants does he want to spend
the rest of his (eternal) life hearing angels play the harp, or is he more interested in playing
roulette with the devils?
Solution of exercise 21.8. This exercise is a simple application of the Bienaym-Tchebychev
inequality. The goal is that the probability of the grade being too far from m be less than 0.2.
Consider the random variable X1 + + X N . It has expectation N m, and S N has expectation
m. The variance of X1 + + X N , on the other hand, is equal to N 2 , as can be seen by looking
at the centered random variables Xi = Xi m, with zero expectation and standard deviation .
The variance of the sum X1 + X2 is simply E ((X 1 + X 2 )2 ) = E ((X1 )2 ) + E ((X2 )2 ) = 2 2 ,
and similarly by induction for the sum of N random variables (since every crossed
term has
p
expectation equal to zero). Hence p
the standard deviation of X1 + + X N is N and that
of S N = (X1 + + Xn )/N is / N .
We are looking for a value of N such that
1
P |S N m|
0.2.
2
The Bienaym-Tchebychev inequality states that we always have

P |S N m| p
2 .
N
p
So we need 1/2 = 0, 2, or = 5, and it then suffices that
p
5
1
p = ,
2
N
or N = 4 5 2 = 20 2 to be certain that the desired inequality holds.
For a value = 2, one finds N = 80!!! This is rather unrealistic, considering the average
budget of a university. If the graders manage to grade very consistently, so that their individual
standard deviation is = 12 , it will be enough to take N = 5. This is still a lot but will
provoke fewer cardiac incidents among university administrators.
However, notice that those estimates are not optimal, because the Bienaym-Tchebychev inequality can be refined significantly if the distribution of the random variables Xi is known (it
is only in itself a worst-case scenario). We know that 80 graders will achieve the desired goal,
independently of the distribution (given the expectation and standard deviation), but it is very
likely that a smaller number suffices.
Solution of exercise 21.9. The probability densities of the masses of the first and second
egg, respectively, are
0
f1 (x) = 1
10
if x
/ [50, 60] ,
if x [50, 60] ,
and
0
f2 (x) = 1
if x
/ [40, 45] ,
if x [40, 45] .
567
The expectation for the mass of the first egg is of course 55 g, and the expectation for the mass
of the second is 42.5 g. Thus the expectation for the total mass is 97.5 g.
The probability density of the mass of the two eggs together is the convolution
Z +
f (x) = f1 f2 (x) =
f1 (t) f2 (x t) dt,
(because, obviously, we can assume that the two masses are independent), which is the function
with the following graph:
0,1
90
95
100 105
The expectation may also be computed from this formula: it is equal to

Z +
E=
x f (x) dx = 97.5,
0
as it should be. The standard deviation satisfies

Z +
2 =
(x E )2 f (x) dx = 12 + 22 .
2
Since 12 = 25/3 and 22 = 25/6
p by immediate computations, we obtain = 25/2. So the
standard deviation is = 5/ 2. From the graph, the probability that the mass is at least 100
g is 1/4.
Solution of exercise 21.11. Let X1 and X2 be two normal centered random variables, with
variances 12 and 22 , respectively. By Proposition 20.67, the probability density of the ratio
X1 /X2 is given by
Z

2 2
t y
1
t2
exp 2 exp 2 |t| dt
g( y) =
21 2
21
22
Z

2

2
1
t
2
2
=
exp 2
+
y
|t| dt
21 2
22 12
1
1
=
,
1 2 22 /12 + y 2
which is the distribution function of a Cauchy
R random variable. (To obtain the last line, notice
that the integral to compute is of the type u(t) u (t) dt.)
Solution of exercise 21.15. The characteristic function of a Cauchy r.v. is given by

f (u) =
a e i ux
dx = e a|u| .
a2 + x 2
Since the characteristic function of the sum of independent random variables is the product
of the characteristic functions of the arguments, we get that the characteristic function of
X1 + + Xn is

exp (a1 + + an ) |x| ,
which does not converge pointwise to the characteristic function of a normal random variable.
568
i) From the independence of the r.v. Xk , we deduce that

Gn =
n
Y
k=1
GXk = (GX )n .
ii) Using the fact that

P(S = n) =
P(S = n and N = k) =
k=0
k=0
we derive
GS (t) =
=
=
P(S = n) t n =
n=0
P
k=0 n=0
k=0
n=0 k=0
P(Sk = n and N = k) t n
P(Sk = n) P(N = k) t n =
P(N = k) Gk (t)
n
P(Sk = n and N = k),
k=0
P(N = k)

= G N GX (t) .
n=0
P(Sk = n) t n
{z
}
Gk (t)

iii) E(S ) = GS (1) = G N GX (1) GX (1) = E(N ) E(X ). An easy computation yields
Var(S ) = GS (1) + GS (1) GS (1)2 = E2 (X ) V (N ) + E(N ) V (X ).
iv) If p0 = 1, we have G 1. If p0 < 1, G is a sum of increasing functions, one of which

at least is strictly increasing.
v) If p0 + p1 = 1, then G is affine, with G (0) > 0 and G (1) = 1. If p0 + p1 < 1, G is the
sum of convex functions, one of which at least is strictly convex, so it is strictly convex.
vi) The graph of G may be of one of the following types:
a) :
b) :
c) :
The equation G (x) = x has two, one, and two solutions, respectively, in the three cases.
vii) Using question ii), we find that G Z2 = G G and then, by induction, that
G Zn+1 = G Zn G = G
G}.
| {z
n+1 times
viii) We have xn+1 = G G G (0) = G (xn ). In all cases above, there exists an interval
[0, ] on which G (x) x, and this interval is stable by G , so that the sequence (xn )nN
is increasing and bounded from above by . Hence, it converges.

ix) We have E(Zn+1 ) = G Z n+1 (1) = G Z n G (1) G (1) = E(Zn ) E(Z ) since G (1) = 1 and
hence, by induction again, we get E(Zn ) = E(Z )n .
x)
(a) There are two cases.

If E(Z) < 1: the curve is as in Figure b) above, which shows that xn 1: there
is certain extinction, and E(Zk ) 0.
If E(Z) = 1 and p0 > 0: the curve is again of type b), so xn 1: there is again
certain extinction and of course E(Zn ) = 1.
569
(b) We assume now that E(Z1 ) > 1. The graph of G is of type c), so that xn
]0, 1]. The probability of extinction is not 1 (there is positive probability because
the particles may all disappear during the first steps). Moreover, E(Zn ) +:
the reaction is explosive.
i) We compute the expectation of Xn :
n
Z 2
n
n
X
X
X
1
cos d = 0,
E(Xn ) = E
cos k =
E(cos k ) =
2 0
k=1
k=1
k=1
so Xn is centered, and of course so is Yn .

There only remains to compute 2 (Xn ) = E(Xn2 ) E(Xn )2 = E(Xn2 ), which is given
by
(
!
2 )
n
n
X
X
X
cos k
cos2 k +
2 (Xn ) = E
=E
cos k cos l
k=1
k=1
k6= l
= n E(cos2 ) + n(n 1) E(cos k cos l ).
Since the variables k are independent, we have

E(cos k cos l ) = E(cos k ) E(cos l ) = 0,
so we obtain finally
2 (Xn ) = n E(cos2 ) =
Similarly, we have
2 (Yn )
n
2
cos2 d =
n
.
2
= n/2.
ii) To determine whether the variables Xn and Yn are correlated or not, we compute the
covariance. We have

Cov(Xn , Yn ) = E Xn E(Xn ) Yn E(Yn ) = E(Xn Yn )
n n
XX
=E
cos k sin l
k=1 l=1
= n E(cos sin ) + n(n 1) E(cos sin )

= n E(cos sin ) + n(n 1) E(cos ) E(sin ) = 0.
Hence Xn and Yn are uncorrelated.

However, they are certainly not independent. For instance, if Xn = n, we necessarily
have Yn = 0 (Pierre has walked due east during the first n steps), so
P(Xn = n and Y = 0) = P(Xn = n) 6= P(Xn = n) P(Yn = 0).
iii) The expectation of Rn2 is
E(Rn2 ) = E(Xn2 + Yn2 ) = E(Xn2 ) + E(Yn2 ) = n.
So the average of the square of thepdistance after n steps is n. Of course, one can not
deduce that the average distance is n (which is in fact false).
iv) Since the r.v. cos n are independent, have zero expectation and finite variance, we may
use the central limit theorem to deduce that the distribution of the r.v. Xn is close to
570

that of a normal centered r.v. with the same variance n/2. The density probability for
such a random variable Xn is
2
x
1
fn (x) = p
exp
.
n
n
Of course, the same is true for Yn , which can be approximated by Yn which also has
probability density fn .
Since Xn and Yn are uncorrelated, it is natural to assume that Xn and Yn are also
uncorrelated, and therefore they are independent, so that the joint probability density
is
2

x + y2
1
exp
.
f Xn ,Yn (x, y)
n
n
v) The distribution function Hn (r ) of the r.v. Rn is given by
ZZZ
Hn (r ) = P(R r ) =
f Xn ,Yn (x, y) dx dy,
B(0 ; r)
where B(0 ; r ) is the closed disc of radius r centered at the origin. Using the approximation above for the density and polar coordinates, we obtain
Z 2 Z r
Z
2 r 2 /n
1
2 /n
e
d d =
e
d
Hn (r )
n 0
n 0
0
Z r2
du
2
=
e u/n
= 1 e r /n .
n
0
The probability density of Rn is therefore
hn (r ) = Hn (r )
2r r 2 /n
e
.
n
vi) Finally, the expectation of Rn is computed by

p
Z +
Z
2 + 2 r 2 /n
n
E(Rn ) =
r hn (r ) dr
r e
dr =
.
n 0
2
0
Notice that the distance to the origin increases with time (rather slowly, compared
with n).
The following heuristic reasoning suggests an explanation: the probability of getting
farther from the pub is larger than the probability of getting closer to the pub, as can
be seen in the picture below: the circular arc describing the values of leading to
an increase of the distance is larger than the other (in fact, a simple computation
p
shows that the difference between the length of the two arcs if of order 1/ r , which is
consistent with the result obtained).
going farther
pub
going closer
What a beautiful argument! Unfortunately, it comes somewhat after the fact, and it
has limits. Indeed, a one-dimensional random walk is also characterized by a distance
571
p
to the origin E(Rn ) proportional to n, whereas the probability of getting closer is the
same as the probability of going further from the origin.
The conclusion to draw from this tale: beware nice heuristic reasonings; sometimes,
they are but webs of lies and deceit.
Appendix
Reminders concerning
topology and normed
vector spaces
A.1
Topology, topological spaces
Since topology is probably not well known to many physicists, we recall a
few elementary definitions. At the very least, they are very useful for stating
precisely properties of continuity and convergence in sets which are more
complicated than R.
DEFINITION A.1 (Open sets, neighborhoods) Let E be an arbitrary set. A
topology on E is the data of a set O of subsets of E , called the open subsets

of E (for the given topology), which satisfy the following three properties:
i) O and E O;
ii) for any set I , and any family (Ai ) iI of open subsets, the union
is also an open set (any union of open sets is open);
iI A i
Topology and normed vector spaces
574
n
iii) for any n N and
Tn any finite family (A1 , . . . , An ) O of open sets,
the intersection i=1 Ai O is open (any finite intersection of open sets is
open).
Finally, a neighborhood V of a point a E is any subset V E

containing an open set V such that a .
If a set E is given a topology O, then the pair (E , O), or E itself by abuse
of notation, is called a topological space.
Remark A.2 If E is a metric space, that is, a set on which a distance function is defined, the
distance can be used to define a metric topology, using the open balls defined by the distance. In
particular this is the case for normed vector spaces, which we will discuss separately below.
Using the notion of open set, the continuity of a map can be defined.
DEFINITION A.3 (Continuity of a map) Let (E , O) and (F , O ) be two topo-
logical spaces. A map f : E

F is continuous if, for any open set F ,
the inverse image f 1 () = x E ; f (x) is an open set in E , tha is,
if f 1 () O for any O .
Remark A.4 This definition should be compared with that of a measurable map between two
measurable spaces (see Definition 2.21 on page 62).
DEFINITION A.5 (Dense subset) Let (E , O) be a topological space. A subset
X E is dense in E if, for any a E and any V which is a neighborhood

of a in E , V meets X , that is, V X 6= .
DEFINITION A.6 (Connected sets) A topological space (E , O) is connected if
it is not the union of two nonempty disjoint open sets: it is not possible to
find A and B such that
X =AB
with A O,
B O,
A 6= ,
B 6= ,
A B = .
A subset X E is connected if (X , O X ) is connected, where the induced

topology O X on X is the topology with open sets of the form U = A X
for some A O.
In other words, the topological space is connected if it is in one piece.
Example A.7 Below are pictures of two sets in the plane; the set (a) is connected, while (b) is
not.
(a)
(b)
Topology, topological spaces
575
Be careful not to mix connected and convex. The first is a purely topological
notion, the second is geometric.
DEFINITION A.8 (Convexity) A subset X in a vector space or in an affine
space is convex if, for any points A , B in X , the line segment [A , B] is
contained in X . (The segment [A , B] is defined as the set of points A +
(1 )B, for [0, 1], and this definition requires a vector space, or affine
space, structure.)
Example A.9 All vector spaces are convex. Below are examples (in the plane) of a convex set (a)
and a nonconvex set (b); note that both are connected.
(a)
(b)
DEFINITION A.10 (Simply connected set) A topological space (E , O) is simply connected if any closed loop inside E may be deformed continuously to
a trivial constant loop (i.e., all loops are homotopic to zero).
A subset X E is simply connected if the topological space (X , O X ) is
simply connected.
In the special case E = R2 (or E = C), a subset X R2 is simply connected
if there is no hole in X .
Example A.11 Still in the plane, the set in (a) is simply connected, whereas the set in (b) is not.
Again both are connected.
(a)
(b)
Example A.12 Let X be the complex plane minus the origin 0; then X is not simply connected.
Indeed, any loop turning once completely around the origin cannot be retracted to a single point
(since 0 is not in X !).
On the other hand, the complex plane minus a half-line, for instance, the half-line of
negative real numbers (i.e., C \ R ) is simply connected. Below are some illustrations showing
how a certain loop is deformed to a point.
Step 1
Step 2
Step 3
Step 4
576
Example A.13 A torus (the surface of a doughnut) is not simply connected. In the picture
below one can see a closed loop which can not be contracted to a single point.
Example A.14 The unit circle U = {z C ; |z| = 1} in the complex plane is not simply
connected.
An important point is that connected and simply-connected sets are preserved by continuous maps.
THEOREM A.15 The image of a connected set (resp. simply connected set) by a con-
tinuous map is connected (resp. simply connected).
DEFINITION A.16 (Compact sets) An open covering of a topological space
K is any family of open sets (U i ) iI such that K is the union of the sets U i :
[
Ui = K .
iI
The topological space K is compact if, from every open covering of K ,

one can extract a finite subcovering, i.e., there exist finitely many indices i1 ,. . . ,
in , such that
in
[
Ui j = K .
j=1
A subset X K is compact if it is compact with the induced topology.

The definition of compact spaces by this property is mostly of interest to
mathematicians. In metric spaces (in particular, in normed vector spaces), this
property is equivalent with the following: from any sequence (xn )nN with values
in K , one can extract a subsequence that converges in K .
Moreover, in any finite-dimensional normed vector space, compact sets are
exactly those which are closed and bounded (see Theorem A.42).
Normed vector spaces
577
A.2
A.2.a Norms, seminorms
In this section, K is either R or C, and E is a K-vector space, not necessarily
of finite dimension.
DEFINITION A.17 (Norm) Let E be a K-vector space. A norm on E is a map
N : E R which satisfies the following properties, for all x, y E and

K:
N1 : we have N (x) 0;

N2 : we have N (x) = 0 x = 0 ;
N3 : we have N ( x) = || N (x);
N4 : we have N (x + y) N (x) + N ( y)
(triangular inequality)
If is customary to write kxk for the norm of an element x, and to denote by

kk the norm map itself.
A normed vector space is a pair (E , N ) where E is a vector space and N
is a norm on E .
DEFINITION A.18 (Seminorm) Let E be a K-vector space. A seminorm on E
is a map N : E R satisfying properties N1, N3, and N4 above. Note that
N3 still implies N (0) = 0.
Example A.19 Let a, b R, a < b . Let C ([a, b ]) be the vector space of continuous functions
defined on [a, b ], with values in K = R or C. For f C ([a, b ]), let

Zb

f (x) dx
k f k1 =
norm of convergence in mean,
a
k f k2 =
sZ

f (x)2 dx
k f k = sup | f |
norm of convergence in mean-square,
norme of uniform convergence.
[a,b]
Then kk1 , kk2 and kk are norms on C ([a, b ]).

On the space of piecewise continuous functions, kk1 and kk2 are seminorms (a function which
is zero almost everywhere satisfies k f k1 = k f k2 = 0 but is not necessarily zero everywhere)
and kk is a norm.
If a, b R, it is necessary to restrict attention to integrable (resp. square integrable, resp.
bounded continuous) functions in order that kk1 , kk2 , or kk be defined; then they are
norms.
578

E , kk be a normed vector space. The
distance associated to the norm kk is the map
DEFINITION A.20 (Distance) Let
d:
E 2 R,
def
(x, y) 7 d (x, y) = kx yk .
Let A E be a non-empty subset of E and let x E . The distance from x

to A is the real number
def
d (x; A) = inf kx yk .
yA
A.2.b Balls and topology associated to the distance

The distance associated to a norm gives a normed vector space the structure of
a metric space. There is then a natural topology associated with this metric; it
is of course a (very) special case of the general topological notions described
in Section A.1.
DEFINITION A.21 Let (E , kk) be a normed vector space. Let a E and let
r R+ . The open ball centered at a with radius r (for the norm kk) is the
set defined by

def
B(a ; r) = x E ; kx ak < r .
The closed ball centered at a with radius r is the set defined by

def
B(a ; r) = x E ; kx ak r .
DEFINITION A.22 (Open set) A subset A E is open in E if and only if,

for any x A, there exists a ball centered at x with positive radius r > 0
contained in A:
x A
r > 0
B(x ; r) A.
DEFINITION A.23 (Closed) A subset A E is closed in E if and only if the

complement of A in E , denoted E A, is open.
Example A.24 A well-known French saying is that Il faut quune porte soit ouverte ou ferme
(Musset; A door must be either open or closed); this is quite true in its way, but not in the
setting of topological spaces and normed vector spaces in particular. For instance, in R2 , the
set ]0, 1[ {0} is neither open nor closed. Similarly, in R, Q is neither open nor closed. More
surprising maybe, R (in R) is both open and closed.
DEFINITION A.25 (Interior) Let A be a subset of E . A point a E is in the
interior of A if there exists r > 0 such that the open ball centered at a with
radius r is contained in A:
r > 0
B(a ; r) A.
The interior of A, denoted A , is the set of all interior points of A.
579
DEFINITION A.26 (Closure) A point a E is in the closure of A if, for

any positive r > 0, the open ball centered at a with radius r has non-empty
intersection with A:
r > 0
B(a ; r) A 6= .
The closure of A, denoted A, is the set of all a satisfying this property.
PROPOSITION A.27 For any A E , we have A A A.

Example A.28 Let A be a non-empty subset of R which has an upper bound. Then sup A A.
PROPOSITION A.29 A subset A E is closed if and only if A is equal to its
closure: A = A. A subset B E is open if and only if B is equal to its interior:
B = B.
The definition of a dense set may then be rephrased in this manner:

DEFINITION A.30 A subset A E is dense in E if and only if the closure
of A is the whole space E .
Example A.31 The set Q of rationals is dense in R; the set of invertible matrices GLn (K) is
dense is the space Mn (K) of all matrices. The set of diagonalizable matrices is dense in Mn (C).
THEOREM A.32 Let E and F be normed vector spaces, D E a dense subset of E ,

and f : E F a continuous map. If f (x) = 0 for all x D , then f is zero on E .
Example A.33 Here is an application: for any matrix A Mn (C), we have det(exp A) =
exp(tr A). This property is obvious in the case of a diagonalizable matrix: indeed, assume that
A = P diag(1 , . . . , n )P 1 (with P GLn (C)); then we have
exp(1 )
n
n
Q
P
..
P 1
exp A = P
and
exp(
)
=
exp
i .
i
.
i=1
i=1
exp(n )
Since the set of diagonalizable matrices is dense in Mn (C), and since the determinant, the trace
and the exponential function for matrices are continuous (the last being not quite obvious), it
follows that the identity is also valid for an arbitrary A Mn (C).
Example A.34 Let f : [0, 1] R be a continuous function such that
x+y
2
f (x) + f ( y)
2
for all (x, y) R2 . We want to prove that f is affine, i.e., there exist a, b R such that
f (x) = a x + b for all x [0, 1]. For this purpose, let b = f (0) and a = f (1) f (0), and let
g : x 7 f (x) (a x + b ). The function g is then continuous, and satisfies g(0) = g(1) = 0.
Moreover, g satisfies the same relation
x+y
g(x) + g( y)
g
=
2
2
580
Edmund Georg Hermann Landau (18771938), German mathematician, was famous in particular for his many works of analytic
number theory (in particular concerning the Riemann zeta function) and his treatise on number theory. From his youngest days,
he was also interested in mathematical problems and games, and
had published two books of chess problems even before finishing
his thesis. He taught at the University of Berlin from 1899 to 1919,
then at the University of Gttingen (replacing Minkowski) until
he was forced to resign by the Nazis in 1934 (the SS mathematician Teichmller organized a boycott of his classes).
as before. From this we deduce immediately that g( 12 ) = 0, then, repeating the process, that
g( 14 ) = g( 24 ) = g( 34 ) = 0,
and by an obvious induction we obtain

k
=0
for all n N and all k [[0, 2n ]].
g
2n

Since the set E = k 2n ; n N , k [[0, 2n ]] is dense in [0, 1], and since g is continuous,
we can conclude that g = 0.
THEOREM A.35 (Properties of open sets) The union of any family of open sets is
an open set.
The intersection of finitely many open sets is an open set.
Example A.36
nN
1
,1
n
= ]0, 1[ is open in R.
Counterexample A.37
The intersection
of infinitely many open sets has no reason to be open:
for instance,
nN
0, 1 +
1
n
= ]0, 1] is not open in R.
THEOREM A.38 (Properties of closed sets) The union of finitely many closed
sets is a closed set.
The intersection of any family of closed sets is a closed set.
Counterexample A.39 As before,
the union of infinitely many closed sets has no reason to be
S
closed: for instance,
nN
0, 1
1
n
= [0, 1[ is not closed.
A.2.c Comparison of sequences
DEFINITION A.40 Let E be a normed vector space.
A sequence (un )nN with values in E is dominated by a sequence (n )nN

of non-negative real numbers, denoted u n = O( n ) and pronounced un is
big-Oh of n , if there exist a real number A and an integer N such that
kun k A |n | for all n N .
581
The Czech priest Bernhard Bolzano (17811848), was born and

died in Prague, in a German-speaking family of Italian origin. He
studied mathematics during his free time. Because of his criticism
of the Austro-Hungarian Empire, he was forbidden to publish
his work and was closely watched. He sought to introduce rigor
in the foundations of mathematics: function theory, set theory,
logic. His impressive works (on cardinal theory, continuity, real
analysis) were in advance on their time and were neglected during
his lifetime (his discovery of the Cauchy criterion for convergence
of sequences was not known). Finally, as a philosopher of science,
he exerted an important influence (notably on Husserl), and he
founded epistemology on formal logic.
The sequence (un ) is negligible compared to (n )nN if, for any positive
real number > 0, there exists an index N N, depending on , such that
kun k |n | for all n N . This is denoted u n = o( n ), pronounced un is
little-Oh of n .
Two real- or complex-valued sequences u = (un )nN and v = (vn )nN are
equivalent, denoted u n v n , if un vn = o(un ) or, equivalently (and it
shows that the relation is an equivalence relation), if un vn = o(vn ).
The notation above is due to E. Landau.
A.2.d Bolzano-Weierstrass theorems

THEOREM A.41 (Bolzano-Weierstrass) Any bounded complex- or real-valued se-
quence has a convergent subsequence.
THEOREM A.42 (Bolzano-Weierstrass) Any subset of a finite-dimensional normed
vector space which is bounded and closed is compact, and in particular, any sequence
with values in such a subset has a convergent subsequence.
A.2.e Comparison of norms

THEOREM A.43 Let E be a vector space and let N , N be two norms on E . Then
any sequence converging to 0 in the sense of N also converges to 0 in the sense of N

if and only if there exists a real number 0 such that
N N ,
that is, N (x) N (x) for all x E .

DEFINITION A.44 Two norms N and N on E are equivalent if there exist
real numbers , 0 such that
N N N .
582
Karl Theodor Wilhelm Weierstrass (18151897), German mathematician (from Westphalia), a famous teacher, was the first to give a convincing construction of the set R of real numbers (such a construction was lacking in Bolzanos work). He also constructed an example
of a continuous function on R which is nowhere differentiable (ignorant of Bolzanos prior work). The famous theorem stating that
any continuous function on an interval [a, b ] can be uniformly approximated by polynomials is also due to him. Finally, wishing to
provide rigorous and unambiguous definitions of the concepts of
analysis, he introduced the definition of continuity based on epsilons and deltas, which are the cause of such happy moments and
memories in the lives of students everywhere.
THEOREM A.45 If the norms N and N are equivalent, then any sequence that
converges for N also converges for N and conversely, and moreover, the limits are
equal.
Example A.46 Let E = Kn . We can easily compare the norms N1 , N2 , and N defined by
N1 ( x) =
|xi |, N2 ( x) =
|xi |2
1/2
and N ( x) = max |xi |. Indeed, we have
N N1 n N
1
hence p N2 N1 n N2 .
p
n
N N2 n N
Thus, those three norms are equivalent.
The previous example illustrates an important theorem in finite-dimensional vector spaces:

Let E be a finite-dimensional vector
space. Then all norms on E are equivalent. In particular, all notions defined in
terms of the topology associated to a norm on E are identical whatever the norm used
in the definition (open sets, closed sets, compact sets, convergence of sequences, limits of
sequences, boundedness,1 Cauchy sequences, Lipschitz functions,..).
THEOREM A.47 (Equivalence of norms)
Counterexample A.48 This is false

for infinite-dimensional vector spaces. For instance, consider the space E = C [0, 1] , R of real-valued continuous functions on [0, 1], with the two
norms
Z1

f (t) dt.
N1 ( f ) = sup f (t)
and
N2 ( f ) =
t[0,1]
It is clear that N2 N1 , so any sequence of functions that converges uniformly to 0 also

converges in mean to 0. However, the converse is false: for instance, let ( fn )nN be the sequence
of functions defined by fn (x) = x n . Then ( fn )nN converges to 0 in means (N2 ( fn ) = 1/(n+1)),
whereas k fn k = 1 for all n N.
Counterexample A.49 On R, the norms of convergence in mean, in quadratic mean, and of
uniform convergence, are pairwise non-equivalent. To see this, let be the function defined
by (x) = 1 |x| if x [1, 1] and (x) = 0 otherwise.
1
The property of being bounded; of course, a bound for a set may depend on the norm.
Exercise
583
The sequence of continuous functions on R defined by fn (x) = (x/n) /n converges

uniformly and in quadratic mean to 0, but it does not converge in mean to 0.
p
The sequence of continuous functions on R defined by gn (x) = n (nx) converges
in mean to 0, but does not converge in quadratic mean, or uniformly.
p
The sequence of continuous functions on R defined by hn (x) = (x/n) / n converges
uniformly to 0, but converges neither in mean nor in quadratic mean.
The reader will now easily find examples illustrating the other cases, using the same ideas.
Another fundamental result specific to finite-dimensional vector spaces is

the following:
THEOREM A.50 (Finite-dimensional normed vector spaces are complete)
The vector spaces R, C, and more generally all finite-dimensional vector spaces are
complete, that is, any Cauchy sequence in a finite-dimensional vector space is convergent.
A.2.f Norm of a linear map

The continuity of a linear map between normed vector spaces is characterized
by the following theorem:
THEOREM A.51 Let E and F be normed vector spaces and let f L(E , F ) be
a linear map from E to F . Then f is continuous if and only if there exists a real
number k R such that k f (x)k F k kxk E for all x E .
If f L(E , F ) is continuous, the real number defined
by
f (x)

F
||| f ||| = sup f (x) F = sup
kxk
kxk E =1
xE \{0}
E
is called the norm of f . Hence, by definition, we have

f (x) ||| f ||| kxk
F
for x E .
The map f 7 ||| f ||| is a norm on the vector space of continuous linear maps from
E to F .
EXERCISE
Exercise A.1 Let E , kk
be a normed vector space and let L (E ) denote the vector space

of continuous maps from E to itself
with the associated norm as above. Assume that (E , kk)
is complete. Prove that L (E ), |||||| is also complete.
584
SOLUTION
Solution of exercise A.1. Let (un )nN be a Cauchy sequence in L (E ). For all x E , we
have

u p (x) uq (x) kxk u p uq ,

and this proves that the sequence un (x) n is a Cauchy sequence in E . By assumption, it
converges; we then let
u(x) = lim un (x)
x E .
n
This defines the map u : E E , which a moments thought shows to be linear. Obviously, we
want to prove that u is continuous, and is the limit of (un )nN .
i) The map u is in L (E ).

By the Cauchy property, there exists N N such that u p uq 1 for all p, q N .
Let x E be such that kxk = 1.

For all p, q N , we have u p (x) uq (x) 1. Letting q go to +, we obtain

u p (x) u(x) 1 for all p N . Hence,

u(x) u(x) u N (x) + u N (x) 1 + |||u N ||| .
This proves that

sup u(x) 1 + |||u N |||
kxk=1
and in particular this is finite, which means that u is continuous.

ii) We have lim un = u in the sense of the norm |||||| on L (E ).
n

Let > 0. There exists N N such that u p uq for all p, q N . Let
x E be such that kxk = 1. Then

u p (x) uq (x) ,
p, q N
and letting q go to + again, we obtain

u p (x) u(x)
p N
(we use the fact that the norm on E is itself a continuous map). Since this holds for all
x E such that kxk = 1, we have proved that

u p u .
p N
And of course, since > 0 was arbitrary, this is precisely the definition of the fact that
(un )nN converges to u in L (E ).
Appendix
Elementary reminders of
differential calculus
B.1
Differential of a real-valued function
B.1.a
Functions of one real variable
Let f : R R be a function which is differentiable at a point a R. By

definition, this means that the limit
f (x) f (a)
f (x + h) f (x)
def
= lim
= lim
xa
h0
x a
h
exists. This also means that there exists a unique real number q R such that
f (x + h) = f (x) + q h + o (h).
h0
In fact, we have q = , and this is the derivative of f at a, that is,

f (a) = q = .
An equivalent rephrasing is that there exists a unique linear map
L (R) such that
f (x + h) = f (x) + .h + o (h),
h0
586
Elementary reminders of differential calculus
since a linear map of R into itself is simply multiplication by a constant

value. (We use the notation .h for the value of evaluated at h, instead
of the more usual (h); this is customary in differential calculus.) This new
formulation has the advantage of generalizing easily to a function of many
variables, that is, to functions defined on Rn .
B.1.b Differential of a function f : R n R
Let f : Rn R be a real-valued function of n real variables, or one vector

variable. We want to generalize the notion of derivative of f , at a point
a Rn . First, the i-th partial derivative of f at the point a Rn is defined
as the limit
f
f ( a + h e i ) f ( a)
def
( a) = lim
h0
xi
h
(if it exists), where e i is the i-th vector of the canonical basis of Rn . A classical
result of differential calculus states that if all partial derivatives f / x i exist
and are continuous functions of a, then we have
f ( a + h) = f ( a) +
f
( a) h i + o ( h).
h0
xi
(b.1)
for h Rn . Denote by dx i the linear form i-th coordinate on Rn , that is,

the linear map associating dx i . h = h i to a vector h = (h1 , . . . , hn ). Then we
can state the following definition:
DEFINITION B.1 (Differential) Let f : Rn R be a map which has continu-
ous partial derivatives at a Rn . The differential of f at a, denoted d f a ,

is the linear map d f a : Rn R defined by
n
X
f
d fa.h =
( a) h i ,
xi
i=1
or equivalently
d fa =
n
X
f
( a) dx i .
xi
i=1
We can then write (b.1) as follows:

f ( a + h) = f ( a) + d f a . h + o ( h).
h0
This turns out to be the right way of defining differentiable functions: f

is differentiable at a if and only if there exists a linear map : Rn R
such that
f ( a + h) = f ( a) + . h + o ( h),
h0
and the linear map is the same as the differential defined above using
partial derivatives.
Differential of map with values in R p
587
B.1.c Tensor notation

In relativity theory, physicists usually use a single letter to denote a vector, and
use exponents over this letter to indicate the coordinates. Using this tensor
notation (see Chapter 16), we can express the relations above as follows:
f (x + h) = f (x) + h f + o(h)
d f = f dx .
and
def
Writing f, = f , we also have

f (x + h) = f (x) + h f,
and
d f = f, dx .
B.2
Differential of map with values in R p
It is very easy to generalize the previous results for vector-valued functions
of a vector variable. Let thus f : Rn R p be a map with vector values.
All previous notions and results apply to each coordinate of f . We therefore
denote
f1 ( a)
.
.
f ( a) =
for all a Rn .
.
f p ( a)
Then, if each coordinate of f has continuous partial derivatives at a Rn ,
we can write down (b.1) for each component:
f i ( a + h) = f i ( a) +
n
X
fi
( a) h j + o( h),
xj
j=1
i = 1, . . . , p.
These relations can be written in matrix notation as follows:

f1
f1
( a)
( a) h1
f1 ( a + h) f1 ( a)

x1

xn

..
..
= .. + ...
.. + o( h).
.
.
.
.

f ( a + h) f ( a) f p ( a) f p ( a) h
p
p
p
x1
xn
(b.2)
Hence the following definition:
DEFINITION B.2 (Differential) Let f : Rn R p be a map such that each
coordinate has continuous partial derivatives at a Rn . The differential
588
of f at a is the linear map d f a : Rn R p

f1
( a)
h1
x1
. .

df a
.. = ..
f
h p ( a)
n
x1
such that

f1
( a) h1

xn
.
..

.
.. .

fp
( a) h p
xn
The differential of f , denoted d f , is the map

d f : Rn L (Rn , R p ),
a 7 d f a .
Thus, the matrix representing the differential d f a at a is the matrix1 of

partial derivatives of f . The relation (b.2) may also be written in the form
f ( a + h) = f ( a) + d f a . h + o ( h).
h0
Using tensor notation, this becomes

f (x + h) = f (x) + h f + o (h)
h0
and
d f = f dx .
B.3
Lagrange multipliers
Using linear forms, it is easy to establish the validity of the method of
Lagrange multipliers.
We start with an elementary lemma of linear algebra.
LEMMA B.3 Let and 1 , . . . , k be linear forms defined on a vector space E of
dimension n. If the forms (1 , . . . , k ) are linearly independent and if

k
T
i=1
Ker i Ker ,
then is a linear combination of the i .

Exercise B.1 Prove this lemma (see page 591).
Now consider the following problem. We wish to find the extrema of a

real-valued function f (x1 , . . . , xn ) on E = Rn . If is differentiable, we know
that a necessary condition (which is not sufficient) for f to be extremal at a =
(a1 , . . . , an ) is that
1
If p = n, it is called the Jacobian matrix of f .
Lagrange multipliers
d f a = 0,
589
f
( a) = 0 i [[1, n]].
xi
that is,
Now let us make the problem more complicated: we are looking for extrema of f on a (fairly regular) subset of Rn , such as a curve in Rn , or
a surface, or some higher-dimensional analogue (what mathematicians call a
differentiable subvariety of E . Assume that this set is of dimension nk and
is defined by k equations (represented by differentiable functions), namely,
(1)
C (x1 , . . . , xn ) = 0
..
S :
(b.3)
.
C (k) (x1 , . . . , xn ) = 0.
In order that the set S defined by those equations be indeed an (n k)dimensional subvariety, it is necessary that the differentials of the equations
C (1) , . . . , C (k) be linearly independent linear forms at every point in S :
(1)
(k)
dC x , . . . , dC x
are linearly independent for all x S .
We are now interested in finding the points on a subvariety S where f
has an extremum.
THEOREM B.4 (Lagrange multipliers) Let f be a real-valued differentiable function on E = Rn , and let S be a differential subvariety of E defined by equations (b.3). Then, in order for a S to be an extremum for f restricted to S , it is
necessary that there exist real constants 1 , . . . , k such that
(1)
(k)
d f a = 1 dC a + + k dC a .
Proof. Assume that a S is an extremum of f restricted to S , and let U be the
tangent space to S at a. This tangent space is an (n k)-dimensional vector space,
where a U is the origin.
If a is an extremum of f on S , we necessarily have d f a . h = 0 for any vector h
tangent to S (if we move infinitesimally on S from a, the values of f cannot increase
or decrease in any direction). Thus we have U Ker d f a .
On the other hand, U is defined by the intersection of the tangent planes to the
subvarieties with equations C (i) ( x) = 0 for 1 i k (because S is the intersection
T
(i)
of those). Hence we have U = ki=1 Ker dC a , and therefore
k
T
i=1
Ker dC a(i) Ker d f a .
There only remains to apply Lemma B.3.
In order to find the points a S where f may have an extremum,

P
(i)
one solves the system of linear equations d f a =
i dC a , where i are
unknowns, together with the equations expressing that a A . In other
590
words, the system to solve is
k
X
f
C (i)
(
a)
( a) = 0,
i
x1
x1
i=1
..
..
.
.
k
(i)
X
C
f
( a)
i
( a) = 0
xn
xn
i=1
and
C (1) ( a) = 0,
..
.
(k)
C ( a) = 0,
(b.4)
where there are n + k equations, with n + k unknowns, the coordinates of

a = (a1 , . . . , an ) and the auxiliary unknowns 1 , . . . , k , called the Lagrange
multipliers. Those auxiliary unknowns have the advantage of preserving symmetry among all variables. Moreover, they sometimes have physical significance (for instance, temperature in thermodynamics is the inverse of a Lagrange multiplier).
Physicists usually proceed as follows: introduce k multipliers 1 , . . . , k
and construct an auxiliary function
def
F ( x) = f ( x)
k
X
i C (i) ( x).
i=1
Look then for points a Rn giving free extrema of F , that is, where the
n variables x1 , . . . , xn are independent (do not necessarily lie on S ). The
condition that the differential of F vanishes gives the first set of equations
in (b.4), and of course the solutions depend on the additional parameters
i . The values for those are found at the end, using the constraints on the
problem, namely, the second set of equations in (b.4).
Example B.5 Let S be a surface in R3 with equation C (x, y, z) = 0. Let r 0 R3 be a point
not in S , and consider the problem of finding the points of S closest to r 0 , that is, those
minimizing the function f ( x) = k x r 0 k2 , with the constraint C ( x) = 0. The auxiliary
function is F ( x) = k x r 0 k2 C ( x), with differential (at a point a R3 ) given by

dF a = 2 ( a r 0 |) grad C ( a) ,
in other words,
dF a . h = 2( a r 0 ) h grad C ( a) h
for all h R3 .
The free extrema of F are the points a such that ( a r 0 ) is colinear with the gradient of C
at a. Using the condition C ( a) = 0, the values of are found. Since grad C ( a) is a vector
perpendicular to the surface S , the geometric interpretation of the necessary condition is that
points of S closest to r 0 are those for which ( a r 0 ) is perpendicular to the surface at the
point a.
Of course, this is a necessary condition, but not necessarily a sufficient condition (the
distance may be maximal at the points found in this manner, or there may be a saddle point).
Remark B.6 Introducting the constraint multiplied by an additional parameter in the auxiliary
function in order to satisfy this constraint is a trick used in electromagnetism in order to fix
the gauge. In the Lagrangian formulation of electromagnetism, the goal is indeed to minimize
591
the action given by the time integral of the Lagrangian under the constraint that the potential
four-vector has a fixed gauge. A gauge-fixing term, for instance ( A)2 in Lorenz gauge, is
added to the Lagrangian density (see [21, 49]).
SOLUTION
Solution of exercise B.1 on page 588. By linear algebra, there exist additional linear forms
k+1 , . . . , n such that (1 , . . . , n ) is a basis of E , since the given forms are linearly independent. Let (e1 , . . . , en ) be the dual
By the definition of a basis, there are constants
Pbasis.
n
i , 1 i n, such that =
i=1 i i . Indeed, we have i = (e i ) using the formula m (eT
) = m defining the dual basis. Let j = i + 1, . . . , n. The same formula
gives e j 1ik Ker i Ker by assumption, hence j = 0 for j > i. This implies
f Vect(1 , . . . , k ), as desired.
Appendix
Matrices
In this short appendix, we recall the formalism of duality (duality bracket or pairing)
and show how it can be used easily to recover the formulas of change of basis, and to
better understand their origin; this can be considered as a warm-up for the chapter
on the formalism of tensors, which generalizes these results.
C.1
Duality
Let E be a finite-dimensional vector space over the field K = R or C, of
dimension n.
DEFINITION C.1 A linear form on E is a linear map from E to K. The dual
space of E is the space E = L (E , K) of linear forms on E . For E a

linear form and x E a vector, we denote
def
, x = (x).
The map , : E E K is called the duality bracket or duality pairing.

DEFINITION C.2 (Dual basis) Let E = (e1 , . . . , en ) be a basis of E . There is a
unique family E = (e1 , . . . , en ) of linear forms on E such that

i, j [[1, n]]
ei , e j = i j .
is a basis of E , and is called the dual basis of E .

The family EP

For x = ni=1 x i ei , we have x i = ei , x for all i [[1, n]]. The elements
of the dual basis E are also called coordinate forms.
Matrices
594
C.2
Application to matrix representation
We will now explain how the formalism of linear forms can be used to clarify the use of matrices to represent vectors, and how to use them to remember
easily the formulas of change of basis.
C.2.a
Matrix representing a family of vectors
DEFINITION C.3 Let E = (e1 , . . . , en ) be a basis of E and let (x1 , . . . , xq ) be q
vectors in E . The matrix representing the family ( x j )1 j q in the basis E

is the (n q)-matrix M = (a i j )1in such that
1 jq
for all j [[1, q]], the column coefficients a1 j , . . . , an j are

the coefficients of x j expressed in the basis E .
In other words, the vector x j is expressed in the basis E by the linear combination
n
P
xj =
a i j ei .
i=1
a11 a1 j a1q
a
21 a2 j a2q
M = .
..
..
.
.
.
.
an1 an j anq
vector x j
seen in E
Using the dual basis, we can state the formula

a i j = ei (x j ) = ei , x j
for all i [[1, n]] and j [[1, q]]. We can express the definition M =
ei , x j i j symbolically in the form
M = matE (x1 , . . . , xq ) = e | x .
C.2.b Matrix of a linear map

As before, let E be a K vector space with a basis B = (b1 , . . . , bn ); in addition,
let F be another K-vector space with basis C = (c1 , . . . , c p ). We denote C =
(c1 , . . . , c p ) the basis dual to B.
Application to matrix representation
595
Let f L (E , F ) be a linear map from E to F , and let A = matB,C ( f )

be the matrix representing the map f in the basis B and C . Recall that, by
definition, this means that the first column of A contains the coefficients
expressing the image f (b1 ) of the first basis vector of B in the basis C , and so
on: the j-th column contains the coefficients of f (b j ) in the basis C .
It follows that for i [[1, p]] and j [[1, n]], the coefficient Ai j is the i-th
coordinate of f (e j ) in the basis C , that is, we have

Ai j = c i f (b j ) = c i , f (b j ) .
This is denoted also
C.2.c

A = c f b .
()
Change of basis
Now let B = (b1 , . . . , bn ) be another basis of the space E . The change of

basis matrix from B to B , denoted P = Pass(B, B ), is defined as follows:
the columns of P represent the vectors of B expressed in the basis B. In
other words, we have
Pi j = i-th coordinate of b j = b i (b j ).
for all (i, j) [[1, n]].
We denote

P = Pass(B, B ) = b b
Notice that the inverse of the matrix P , which is simply the change of basis
matrix from B to B, is
1

b b
= b b .
C.2.d Change of basis formula
Let again B = (b1 , . . . , bn ) and B = (b1 , . . . , bn ) be two bases of E , and let

x E be a vector. The change of basis formula
matB (x) = Pass(B, B ) matB (x)
expresses the relation between the coordinates for x in each basis. In abbreviated form, we write
X = PX
or, with the previous notation,

b | x = b b b x .
Matrices
596
To remember this formula, it suffices to use the relation

n

P
b b
Id =
k
k
k=1
or its short form

Id = b b
(c.1)
to recover the formula.
Let now f : E F be a linear map. In addition to the basis B and B of
E , let C and C be basis of F , and let
M = matB,C ( f )
M = matB ,C ( f ).
and
If we denote by P and Q respectively the change of basis matrices for E and

F , i.e.,

P = Pass(B, B ) = b b
and
Q = Pass(C , C ) = c c ,
the change of basis formula relating M and M can be recovered by writing

M = c f b = c c c f b b b = Q 1 M P
(c.2)
| {z } | {z } | {z }
M
Q 1
One should remember that (c.1) provides a way to recover without mistake the
change of basis formula for matrices, in the form

c f b = c c c f b b b .
Exercise C.1 The trace of an endomorphism f L (E ) is defined by the formula
tr ( f ) =
i=1
e i , f (e i )
in a basis E = (e1 , . . . , en ). Show that this definition is independent of the chosen basis.
C.2.e
Case of an orthonormal basis
When E has a euclidean or Hilbert space structure and the basis E =

(e1 , . . . , en ) is an orthonormal basis, the coordinate forms ei are associated
with the orthogonal projections on the lines spanned by the vectors ei : we
have
x E i [[1, n]]
ei (x) = (ei |x) .
The formula of the previous section, provided all bases involved are orthonormal, becomes

c f b = c c c f b b b .
or, with all indices spelled out precisely
n P
n

P

c i f (b j ) =
c i ck ck f (b ) b b j .
k=1 =1
Appendix
A few proofs
This appendix collects a few proofs which were too long to insert in the main
text, but are interesting for various reasons; they may be intrinsically beautiful, or
illustrate important mathematical techniques which physicists may not have other
opportunities to discover.
THEOREM 4.27 on page 96 (Winding number) Let be a closed path in C,
U the complement of the image of in the complex plane. The winding number
of around the point z ,
Z
1
d
def
for all z U
(d.1)
Ind (z) =
2i z
is an integer, and as a function of z it is constant on each connected component of U ,

and is equal to zero on the unbounded connected component of U .
Let z . Assume that is parameterized by [0, 1], so that
Z1
1
(t)
Ind (z) =
dt.
2i 0 (t) z
()
If w is a complex number, it is equivalent to show that w is an integer or to show that

e 2i w = 1. Hence, to prove that the winding number Ind (z) is an integer, it suffices to prove
that (1) = 1, where
Z t

(s)
def
(t) = exp
ds
(0 t 1).
0 (s) z
A few proofs
598

Now, we can compute easily that (t)/(t) = (t)/ (t) z , for all t [0, 1], except
possibly at the finitely many points S [0, 1] where the curve is not differentiable. Hence
we find

(t)
(t) (t)
(t)
=
2
(t) z
(t) z
(t) z
=
(t) (t)
(t) (t)
2
2 = 0
(t) z
(t) z

on [0, 1] \ S . Because the set S is finite, it follows that (t)/ (t) z is a constant. In
addition, we have obviously (0) = 1, and hence
(t) =
(t) z
.
(0) z
Finally, since (1) = (0), we obtain (1) = 1 as desired.

It is clear from the definition of Ind as function of z that it is a continuous function on
. But a continuous function defined on which takes integral values is well known to be
constant on each connected component of .
Finally, using the definition () it follows that Ind (z) < 1 if |z| is large enough, and
therefore the integer Ind must be zero on the unbounded component of .
THEOREM 4.31 on page 97 (Cauchy for triangles) Let be a closed triangle
in the plane (by this is meant a filled triangle, with boundary consisting
of three line segments),
entirely contained in an open set . Let p , and let
f H \ { p} be a function continuous on and holomorphic on except
possibly at the point p. Then
Z
f (z) dz = 0.
Let [a, b ], [b , c ], and [c , a] denote the segments which together form the boundary of
the filled triangle . We can assume that a, b , c are not on a line, since otherwise the result is
obvious. We now consider three cases: p
/ , p is a vertex, p is inside the triangle,
but not a vertex.
p is not in . Let c be the middle of [a, b ], and a and b the middle, respectively, of
[b , c ] and [c , a]. We can subdivide into four smaller similar triangles i , i = 1, . . . , 4,
as in the picture, and we have
def
J =
f (z) dz =
4 Z
X
i=1
b
f (z) dz
c
a
c
b
(the integrals on [a , b ], [b , c ], and [c , a ] which occur on the right-hand side do so
in opposite pairs
R because of orientation). It follows that for one index i {1, 2, 3, 4} at

least, we have f (z) dz | J /4|. Let (1) denote one triangle (any will do) for which
i
this inequality holds. Repeating the reasoning above with this triangle, and so on, we
A few proofs
599
obtain a sequence ((n) )nN of triangles such that (1) (2) (n)
and
Z

|J|

.
f
(z)
dz
(n)

4n
at each step. Let L be the length of the boundary . Then it is clear that the length
of (n) is equal to L/2n for n 1.
Finally, because the sequence of triangles is a decreasing (for inclusion) sequence of
compact sets with diameter tending to 0, there exists a unique intersection point:
\
(i) = {z0 }.
i1
By definition, we have z0 and by assumption, f is therefore holomorphic at z0 .

Now we can show that J = 0.
Let > 0. There exists a positive real number r > 0 such that
| f (z) f (z0 ) f (z0 )(z z0 )| |z z0 |
for all z C such that |z z0 | < r . Since the diameter of (n) tends to zero, we have
|z z0 | < r for all z (n) if n is large enough. Using Proposition 4.30, we have
Z
Z
Z
f (z0 ) dz = f (z0 )
dz = 0
and
f (z0 ) (z z0 ) dz = 0.
(n)
(n)
Hence
f (z) dz =
(n)
(n)
which proves that
(n)
[ f (z) f (z0 ) f (z0 ) (z z0 )] dz,

L L
L2
f (z) dz n n = n ,
2 2
4
(n)
Z

J 4n
f (z) dz L2 .
and therefore
(n)
This holds for all > 0, and consequently we have J = 0 in the case where p
/ .
p is a vertex, for instance p = a. Let x [a, b ] and y [a, c ] be two arbitrarily chosen
points. According to the previous case, we have
Z
f (z) dz =
[x yb]
and hence
b
f (z) dz = 0
[ yb c ]
f (z) dz =
p
Z
[a,x]
[x, y]
f (z) dz,
[ y,a]
which tends to 0 as [x a] and [ y a], since f is bounded, and the length of the
path of integration tends to 0.
p is inside the triangle, but is not a vertex. It suffices to use the previous result applied to
the three triangles [a, b , p], [b , p, c ], and [a, p, c ].
A few proofs
600
THEOREM 7.19 on page 187 (Cauchy principal value) Let D(R) be a
test function. The quantity

Z
has a finite limit as [
(x)
dx +
x
+
+
(x)
dx
x
Moreover, the map

Z
(x)
7 lim+
dx
0
|x|> x
0+ ].
(d.2)
is linear and continuous on D(R).
R Let D, so is of C class and has bounded support. In particular the integral

(x)/x
dx
exists
for
all
> 0. Let [M , M ] be an interval containing the support of .

|x|>

For all x [M , M ], we have (x) (0) = x (0) + (x) , where is a function such
that lim x0 (x) = 0.
Z
(0)
Notice that
dx = 0 for all > 0,
x
<|x|<M
so that, by parity, we have

Z
def
f () =
<|x|<M

x (0) + (x)
dx
x
<|x|<M
Z

=
(0) + (x) dx,
(x)
dx =
x
<|x|<M
and as [ 0+ ] this last expression obviously has a limit.

We know must show that the linear map defined for test functions by (d.2) is continuous
on D(R).
Let (n )nN be a sequence of test functions converging to 0. By definition, there exists a
fixed interval [M , M ] containing the support of all the n .
Let > 0. There exists an integer N such that, if n N , we have kn k /M ,
kn k , and kn k . Moreover, the Taylor-Lagrange inequality shows that for n N
and x R, we have
x 2
|n (x) n (0) x n (0)|
k k .
2
Let n N be an integer. Choose, arbitrarily for the moment, a real number ]0, 1], and
consider
Z
Z
Z
n (x)
n (x)
n (x)
dx =
dx +
dx.
x
x
x
|x|>
<|x|<1
|x|>1
The second integral is bounded by 2M kn k 2. Using the symmetry of the interval in
the first integral, we can write
Z
Z
n (x) n (0)
n (x)
dx =
dx.
x
x
<|x|<1
<|x|<1
But we have
Z
Z
Z

n (x) n (0)
n (x) n (0) xn (0)
xn (0)

dx
dx +
dx

<|x|<1
<|x|<1
<|x|<1 x

x
x
Z
x

kn k
dx + kn k 2.
<|x|<1 2
A few proofs
601
This is valid for any ]0, 1]. Letting then [ 0+ ] for fixed n, we find that

Z

n (x)

dx 4
n N .
lim+

0 |x|> x
This now proves that the map (d.2) is continuous.
THEOREM 8.18 on page 232 Any Dirac sequence converges weakly to the Dirac
distribution in D (R).
Let ( fn )nN be a Dirac sequence and let A > 0 be such that fn is non-negative on [A , A]
for all n.
def
Let D be given. We need to show that fn , (0). First, let = (0), so
that is a function of C class, with bounded support, with (0) = 0. Let K be the support
of .
Let > 0. Since is continuous, there exists a real number > 0 such that
|(x)|
for all real numbers x with |x| < . We can of course assume that < A. Because of
Condition in the definition 8.15 on page 231 of a Dirac sequence, there exists an integer N
such that
Z

fn (x) (x) dx

|x|>

for all n N . Because fn is non-negative on [, ], we can also write
Z

Z
Z

fn (x) dx =
fn (x) (x) dx
fn (x) dx.

|x|

|x|
|x|
Again from the definition, the last integral tends to 1, so that if n is large enough we have
Z

Z

fn (x) (x) dx 3.
fn (x) (x) dx 2,
hence

|x|

R
Now, there only remains to notice that K fn (x) dx 1 (since K is a bounded set), and
n
then for all n large enough, we obtain

Z

fn (x) (x) dx (0) 4.

Since is an arbitrary test function, this finally shows that fn .
A few proofs
602
THEOREM 9.28 on page 261 The space 2 , with the hermitian product
(a|b) =
n=0
is a Hilbert space.
a n bn ,
The inequality || 21 (2 + 2 ) shows that 2 is a vector space. It is also immediate to

check that the formula above defines a hermitian product on 2 . The hard part is to prove that
2 is a complete space for the norm kk2 associated to this hermitian product.
Let (u ( p) ) pN be a Cauchy sequence in 2 . We need to prove that it converges in 2 . There are
three steps: first a candidate u is found that should be the limit of (u ( p) ) pN ; then it is proved
that u is an element of 2 ; and finally, the convergence of (u ( p) ) pN to u in 2 is established.
Looking for a limit Since the sequence (u ( p) ) pN is a Cauchy sequence, it is obvious
that for any p N, the sequences of complex numbers (u ( p) n )nN is a Cauchy sequence in C.
Since C is complete, this sequence converges, and we may define
def
un = lim un( p) .
p
Obviously, the sequence (un )nN of complex numbers should be the limit of (u ( p) ) p .
The sequence u is in 2 The Cauchy sequence (u ( p) ) pN , like any other Cauchy sequence, is bounded in 2 : there exists A > 0 such that

X

( p)
u =
u ( p) 2 A
n
2
n=0
for all p N. Let N be an integer N N. We obtain in particular

p N,
and as [p ], we find
N
X

u ( p) 2 u ( p) A ,
n
2
n=0
N
N
2 X
X

|un |2 A
lim un( p) =
n=0
n=0
(since the sum over n is finite, exchanging limit and sum is obviously permitted). Now this
holds for all N N, and it follows that u 2 and in fact kuk
2 A.
The sequence u ( p) tends to u in 2 that is, we have u ( p) u 2 0.

p
Let > 0. There exists an integer N N such that u ( p) u (q) if p, q N . Let
2
p N . Then for all q p and all k N, we have
k
X

u ( p) u (q) 2 u ( p) u (q) .
n
n
2
n=0
We fix k, and let [q ], deriving that
k
X

u ( p) un 2
n
n=0
(as before, the sum is finite, so exchanging limit

sum is possible). Again, this holds for all
and
k N, so in the limit [k ], we obtain u ( p) u 2 .

All this shows that u ( p) u 2 0, and concludes the proof.
p
A few proofs
603
THEOREM 12.36 on page 349
i
1 h
(r c t) (r + c t) ,
4c r
1 (r c t) if t > 0,
F 1 [cos k c t] = 4r
( r)
if t = 0.
F 1
sin k c t
kc
We prove the second formula. It is obviously correct for t = 0, and if t > 0, we compute
ZZZ
1
F 1 [cos k c t] ( r) =
cos(k c t) e i k r d3 k
(2)3
Z 2
Z Z +
1
=
d
d
dk cos(k t) e i kr cos k 2 sin
(2)3 0
0
0
Z +
(e i kr e i kr ) 2
1
cos(k t)
k dk
=
2
(2) 0
ikr
Z +

1
=
e i k(tr) e i k(tr) + e i k(t+r) e i k(t+r) ik dk
2
8 r 0
Z +

1
=
e i k(tr) + e i k(t+r) ik dk
by parity
82 r
h
i

1
1
=
F 1 ik (r
t) + (r
+ t) =
(r t) + (r + t)
4r
4r
by the properties of Fourier transforms of derivatives. Since t > 0, the variable r is positive
and the distribution (r + t) is identically zero, leading to the result stated.
The first formula is obtained in the same manner.
A few proofs
604
THEOREM 16.18 on page 439 (Existence and uniqueness of the tensor
product) Let E and F be two finite-dimensional vector spaces over K = R or C.

There exists a K-vector space E F such that, for any vector space G , the space of
linear maps from E F to G is isomorphic to the space of bilinear maps from
E F to G , that is,
L (E F , G ) Bil(E F , G ).
More precisely, there exists a bilinear map

: E F E F
such that, for any vector space G and any bilinear map f from E F to G , there
exists a unique linear map f from E F into G such that f = f , which is
summarized by the following diagram:
EF
@
@ f
@
@
?
R
@
- G
EF
The space E F is called the tensor product of E and F . It is only unique up to

(unique) isomorphism.
The bilinear map from E F to E F is denoted and also called the
tensor product. It is such that if (e i ) i and (f j ) j are bases of E and F , respectively,
the family (e i f j ) i, j is a basis of E F .
All this means that the tensor product linearizes what was bilinear, or also that
to any bilinear map on E F , one can associate a unique linear map on the
space E F .
Construction of the tensor product First, consider the set E F of pairs (x, y) with
x E , y F . Then consider the free vector space M generated by E F , that is, the set of
abstract linear combinations of the type
X
(x, y) [(x, y)],
(x, y)E F
where [(x, y)] is just some symbol associated with the pair (x, y) and (x, y) is in K, with the
additional condition that (x, y) = 0 except for finitely many pairs (x, y).
The set M has an obvious structure of a vector space (infinite-dimensional), with basis the
symbols [(x, y)] for (x, y) E F : indeed we add up elements of M term by term, and
multiply by scalars in the same way:
X
X
X
(x, y) [(x, y)] +
(x, y) [(x, y)] =
((x, y) + (x, y) )[(x, y)],
(x, y)E F
(x, y)E F
X
(x, y)E F
(x, y) [(x, y)] =
(x, y)E F
(x, y)E F
(x, y) [(x, y)].
A few proofs
605
Now, in order to transform bilinear maps on E F into linear maps, we will use M and
identify in M certain elements (for instance, we want the basis elements [(x, 2 y)] and [(2x, y)]
to be equal, and to be equal to 2[(x, y)]).
For this purpose, let N be the subspace of M generated by all elements of the following
types, which are obviously in M :
[(x + x , y)] [(x, y)] [(x , y)];
[(x, y + y )] [(x, y)] [(x, y )];
[(a x, y)] a[(x, y)];
[(x, a y)] a[(x, y)];
where x E , y F and a K are arbitrary.
Consider now T = M /N , the quotient space, and let be the canonical projection map
M M /N , that is, the map which associates to z M the equivalence class of z, namely, the
set of elements of the form z + n with n N . This map is surjective and linear.
Now, E F can be identified with the canonical basis of M , and we can therefore
construct a map
: E F T = M /N
by sending (x, y) to the class of the basis vector [(x, y)].
We now claim that this map and T as E F satisfy the properties stated for the tensor
product.
Indeed, let us first show that is bilinear. Let x, x E and y F . Then (x + x , y) is the
class of [(x + x , y]). But notice that the element [(x + x , y)] [(x, y)] [(x , y )] is in N by
definition, which means that the class of [(x + x , y)] is the class of [(x, y)] + [(x , y)]. This
means that (x + x , y) = (x, y) + (x , y). Exactly in the same manner, the other elements
we have put in N ensure that all properties required for the bilinearity of are valid.
We now write E F = T , and we simply write (x, y) = x y.
Consider a K-vector space G . If we have a linear map T = E F G , then the composite
E F E F G
is a bilinear map E F G . Conversely, let B : E F G be a bilinear map. We construct
a map E F G as follows. First, let B : M G be the map defined by
X
X
(x, y) [(x, y)] 7
(x, y) B(x, y).
(x, y)E F
(x, y)E F
It is obvious that B is linear. Moreover, it is immediate that we have B(n)

= 0 for all n N
because B is bilinear; for instance,
B([(x
+ x , y)] [(x, y)] [(x , y)]) = B(x + x , y) B(x, y) B(x , y) = 0.
This means that we can unambiguously use B to induce a map T G by x y 7 B(x, y),
because if we change the representative of x y in M , the value of B(x, y) does not change.
Thus we have constructed a linear map E F G .
It is now easy (and an excellent exercise) to check that the applications just described,
L (E F , G ) Bil(E F , G ) and Bil(E F , G ) L (E F , G ),
are reciprocal to each other.
Uniqueness Let : E F H and : E F H be two maps satisfying the

property stated in the theorem, where H and H are K-vector spaces.
Since is itself bilinear, there exists (by the property of ) a unique linear map : H
H such that = . Similarly, exchanging the roles of and , there exists a unique
: H H such that = . Now notice that we have both = IdH (a triviality)
and = . The universal property of again implies that = IdH . Similarly,
we find = IdH , and hence is an isomorphism from H to H .
This shows that the tensor product is unique up to isomorphism, and in fact up to unique
isomorphism.
A few proofs
606
(Poincar formula) Let n N and let

A1 , . . . , An be arbitrary events, not necessarily disjoint. Then we have
!
n
n
[
X
X
P
Ai =
(1)k+1
P(Ai1 Aik ).
THEOREM 19.18 on page 516
1i1 <i2 <<ik n
k=1
i=1
Following Pascal,1 consider the random variables 1Ai which are characteristic functions of
the events Ai , that is,
1Ai : R,
0
7
1
if A ,
if
/ A.
Then we have E (1Ai ) = P(Ai ). Moreover, these random variables satisfy

1Ai A j = 1Ai 1 A j
Using both properties, we find
1 A c = 1 1 A i .
and
P(A1 An ) = E 1
n
Y
i=1
(1 1Ai ) .
Now expand the product on the right-hand side, noting that 1 1Ai = 1Ai ; we get2
n
n
Y
X
(1 1Ai ) = 1 +
(1)k
i=1
k=1
1i1 <i2 <<ik n
1 Ai 1 Ai . . . 1 Ai .
1
It only remains to compute the expectation of both sides to derive the stated formula.
Blaise Pascal (16231662) studied mathematics when he was very young, and investigated
the laws of hydrostatics and hydrodynamics (having his brother in law climb to the top of
the Puy de Dme with a barometer in order to prove that atmospheric pressure decreases with
altitude). He found many results in geometry, arithmetics, and infinitesimal calculus, and was
one of the very first to study probability theory. In addition, he invented and built the first
mechanical calculator. In 1654, Pascal abandoned mathematics and science for religion. In his
philosophy there remain, however, some traces of his scientific mind.
2
This may be proved by induction, for instance, in the same manner that one shows the
that
n
n
Y
X
X
(1 xi ) = 1 +
(1)k
xi1 xi2 . . . xik .
i=1
k=1
1i1 <i2 <<ik n
Tables
609
Fourier Transforms
The convention used here is the following:
Z +
Z
def
2i
x
e
f () =
f (x) e
dx
and
f (x) =
For other standard conventions, see the table on page 612.

fe()
f (x)
f (a x)
f (n) (x)
(2ix)n f (x)
(x)
x
b a [a,b]
e |x|
a2
2a
+ 42 2
e x
1
x
xk
fe() e 2i x d.
1 e
f
|a|
a
(2i)n fe()
fe(n) ()
sin
= sinc
sin |a|

sin (b a) i(a+b )
e
2
2
+ 42 2
e a||
r
2 2 /
e
2i
1
(k)
(2i)k
(b > a)
( > 0)
(a > 0)
( > 0)
Tables
610
f (x)
H (x)
H (x)
sgn(x)

1
pv
x
sin(20 x)
cos(20 x)
(x)
fe()
1
1
+
pv
2 2i
1
1
pv
2 2i

1
1
pv
i
i sgn

1
( 0 ) ( + 0 )
2i

1
( 0 ) + ( + 0 )
2
1
(x x0 )
e 2i x0
X(x)
X()
1
X(/a)
|a|
X(a x)
Fourier Transforms
611
Three-dimensional Fourier transforms

For the following three-dimensional Fourier transforms, the convention
chosen is
ZZZ
ZZZ
1
F ( k) =
f ( x) e i k r d3 r ,
f ( r) =
F ( k) e i k r d3 k.
(2)3
Moreover, we put r = k rk, n = r/r and k = k kk.
f ( r)
F ( k)
f ( r)
k 2 F ( k)
1/r
4/k 2
e r /r
1
2
r + a2
4/(2 + k 2 )
( > 0)
e ak /22 k
(a > 0)
1/r 2
1/22 k
1
( p n)( q n)
r
e r
( r)
1
4
k2
( p k)( q k)
p q 2
k2
3/2
2
e k /4
(2)3
p, q R3
( > 0)
f (x) e 2i x dx
fe( p) =
f ( r ) e i k r d3 r
f ( x) e i p x/ h}
x
(2 h})3/2
d3
dx
f (x) e i p x/ h} p
2 h}
ZZZ
ZZZ
fe(p) =
F ( k) =
fe() e 2i x d
Z
1
f (x) = p
fe() e i x d
2
Z
1
f (x) =
fe() e i x d
2
ZZZ
1
f ( r) =
F ( k) e i kr d3 k
(2)3
Z
dp
f (x) = fe(p) e i p x/ h} p
2 h}
ZZZ
d3 p
f ( x) =
fe( p) e i p x/ h}
(2 h})3/2
f (x) =
Inverse transform
(2 h})3/2
p
2 h}
2 h}
(2 h})3/2
2
(2)3
1
p
2
F []
F [1]
f g=
f g=
fe e
g
fe e
g
fe e
g
f g=
f g=
fe e
g
Z
1
fe e
g
(2)3
Z
Z
f g = fe e
g
1
f g=
2
Parseval-Plancherel
Various definitions used for the Fourier transforms. The formulas for those different from the ones used in this book can be recovered by elementary changes
of variables. The convention named QM are used in quantum mechanics. Note that the simplest definition in terms of memorizing the formulas is the
first one. The definition for quantum mechanics is the one in the book by Cohen-Tannoudji et al. [20]
(3D)
QM
(1D)
QM
3D
Z
1
fe() = p
f (x) e i x dx
2
Z
e
f () = f (x) e i x dx
fe() =
Fourier transform
612
Tables
Table of the usual conventions for the Fourier transform
Laplace transforms
613
Laplace transforms
In some problems (in particular, to solve differential equations while taking initial conditions at t = 0 into account), it is useful to go through the Laplace transform, which transforms
the differential operations into algebraic operations, which can easily be inverted. The table
of Laplace transforms is therefore usually read from right to left: one looks for an original
t 7 f (t) for the function p 7 F (p) that has been calculated.
f (t)
1
2i
c +i
F (p) =
e t z F (z) dz
e p t f (t) dt
0
F (p)
c i
f (t)
inversion
p F (p) f (0+ )
F (p)
p
f (s) ds
(1)n t n f (t)
F (n) (p)
Z
f (t)
t
e at f (t)
F (z) dz
(n N)
()
F (p a)
(a C)
R
() The integral p is defined on any path joining p to infinity on the right-hand side
(Re(z) +), which is entirely contained in the half-plane of integrability where F is
defined. Because F is analytic, this integral is independent of the chosen path.
Tables
614
f (t)
F (p)
(t)
(n) (t)
pn
X(t)
1
1 ep
1/p
t
p
1/p 2
p
1/ p
e at
1/(p a)
(a C)
t n e at /n!
(p a)(n+1)
(n N,
1/ t
cos t
p2
p
+ 2
sin t
p 2 + 2
cosh t
p
p 2 2
sinh t
p 2 2
1
(1 cos t)
2
1
p(p 2 + 2 )
1
(t sin t)
3
1
p 2 (p 2 + 2 )
1
(sin t t cos t)
23
1
(p 2 + 2 )2
t
cos t
2
1
(sin t + t cos t)
2
(p 2
p
+ 2 )2
p2
(p 2 + 2 )2
(n N)
a C)
( R or C)
( C)
Laplace transforms
f (t)
F (p)
t cos t
p 2 2
(p 2 + 2 )2
J0 (t)
J1 (t)
J1 (t)
t
log t
p
1
t k 2
Jk 1 (t)
2
(k) 2
1
p2
1 p
p
615
+ 2
p2
+ 2
p 2 + 2 p
log p
1
(p 2 + a 2 )k
( 0, 577...)
(k > 0)
[a, b]
R
R+
R
R+
Uniform law
Normal law
Exponential law
Cauchy law
2 law
a
1
2
a + (x m)2
x n 1
1
2
e x/2
2(n/2) 2
1 e x/
not defined
2n
not def.
(1 2i)n/2
e im|a|
1
1 i
2 2 /2
2
m
e im e
e ib e ia
i(b a)
(b a)2
12
a+b
2
b a [a,b ]

1
(x m)2
p exp
2 2
2
Char. function ()
Continuous laws
n p(1 p)
Variance 2
Poisson law
np
n2 1
12
p(1 p)
n+1
2
p
1
P(X = k) =
n
P(X = 1) = p
P(X =0)= 1 p
n k
P(X = k) =
p (1 p)nk
k
n
P(X = n) = e
n!
Expectation E (X )
[[0, n]]
Binomial law B(n, p)
Variance
Expectation
Law
Density f (x)
{0, 1}
Bernoulli law B(1, p)
Set of values
[[1, n]]
Uniform law
Name
Set of values
Name
Discrete laws
616
Tables
Tables of probability laws
Further reading
General books
Methods of theoretical physics
P. M. Morse and H. Feshbach (McGraw-Hill, 1953)
A great classic, very complete.
Mathematical methods for physicists

G. B. Arfken and H. J. Weber (Academic Press, 4th edition, 1995)
Many subjects are covered, in particular, differential equations, Bessel and Legendre functions, other special functions, spherical harmonics.
Integration
Lebesgue measure and integration, an introduction
F. Burk (Wiley Interscience, 1998)
A very clear book, with many examples guiding the reader step by step. Also a
very good historical introduction.
Complex analysis
Function theory of one complex variable
R. E. Greene and S. G. Krantz (Wiley Interscience, 1997)
Very pedagogical.
Distributions and Fourier analysis

Fourier analysis: an introduction
E. M. Steinz and R. Shakarchi (Princeton University Press, 2003)
A very clear introduction.
A wavelet tour of signal processing

S. Mallat (Academic Press, 1999)
A nice and clear treatment of a difficult subject, with many illustrations and
examples of practical applications of wavelets.
Bibliography
618
Integral transforms, operator theory, Hilbert spaces

Theory of linear operators in Hilbert space
N. I. Akhiezer and I. M. Glazman (2 volumes, Dover, 1993)
A very complete book, rather difficult to read.
Applied functional analysis

E. Zeidler (Springer-Verlag, 1995, Applied Mathematical Sciences 108)
This book is oriented toward physical applications, but remains quite difficult.
Quantum mechanics in Hilbert space

E. Prugovecki (Academic Press, 1971, Pure and Applied Mathematics 41)
Another very complete treatement, mathematically rigorous, using powerful
tools which go far beyond what we have seen in the present text.
Tensors, geometry
Riemannian geometry
S. Gallot, D. Hulin and J. Lafontaine (Springer-Verlag, 1990)
Leon sur la thorie des spineurs
. Cartan (Hermann, 1966)
An English version is published by Dover with the title The theory of spinors, 1981.
Introduction to differentiable manifolds

L. Auslander and R. E. MacKenzie (Dover, 1977)
A clear introduction to differential geometry.
Modern geometryMethods and applications

B. Dubrovin, S. Novikov and A. Fomenko (Springer, 1991)
Three volumes: I. Geometry of surfaces, of transformation groups and fields; II. Geometry
and topology of varieties; III. Methods of homology theory. The first two volumes
are most likely to be of interest to physicists. The geometrical intuition is
emphasized throughout. The English edition published by Springer-Verlag is
much more expensive than the original Russian or French editions.
Geometry, topology and physics

M. Nakahara (Adam Hilger, 1990)
A unique book that makes the link between theoretical physics and high-level
mathematics.
Groups
Symmetries in quantum mechanics: from angular momentum to supersymmetry
M. Chaichian and R. Hagedorn (Institute of Physics, 1998)
A very interesting book, with increasing difficulty, giving a vast overview of symmetries, and explaining clearly the relation between symmetries and dynamics.
Bibliography
619
Lie groups and algebra, with applications to physics, geometry and mechanics
D. H. Sattinger and O. L. Weaver (Springer-Verlag, 1986, Applied Mathematical Sciences 61)
Probabilities, statistics
An introduction to probability theory
W. Feller (Wiley, 1968 (vol. I) and 1971 (vol. II))
The first volume, with little formalism, is a remarkable book directed toward
concrete applications, analyzed in great detail and depth.
Stochastik
A. Engel (Ernst Klett Verlag, 2000)
This book has a fascinating approach to probability theory, full of applications
and with little formalism. Moreover, this is one of the rare books dealing with
probabilities and statistics simultaneously. (A French translation of the original
German edition is available: Les certitudes du hasard, ALAS diteur, Lyon, 1990).
References
[1] Edwin A. Abbott. Flatland. Dover, 1992.
[2] Milton Abramowitz and Irene A. Stegun, editors. Handbook of mathematical functions. Dover, 1972.
[3] Sir George Airy. On the intensity of light in the neighbourhood of a
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[4] Naum Ilitch Akhiezer and Isral Markovitch Glazman. Theory of linear
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[15] Henri Broch. Au cur de lextraordinaire. Collection Zttique, 2002.

A book about the spirit of doubt a priori, which should be always
present to the mind of scientists...
[16] Frank Burk. Lebesgue measure and integration. Wiley-Interscience, 1998.
[17] Henri Cartan. Cours de calcul diffrentiel. Hermann, 1979. (Differential
forms, Dover, 2006).
[18] Masud Chaichian and Rolf Hagedorn. Symmetries in quantum mechanics. Institute of Physics Publishing, 1998.
[19] ric Charpentier, Annick Lesne, and Nikola Nikolski. Lhritage de
Kolmogorov en mathmatiques. Belin, 2004.
[20] Claude Cohen-Tannoudji, Bernard Diu, and Frank Lalo. Quantum
Mechanics. WileyInterscience, 2001.
[21] Claude Cohen-Tannoudji, Jacques Dupont-Roc, and Gilbert Grynberg. Photons and atoms, Introduction to quantum electrodynamics. Wiley
Interscience, 1997.
[22] Jean-Franois Colombeau. Elementary introduction to new generalized functions. North Holland, 1984.
[23] Jean-Franois Colombeau. New generalized functions and multiplication of
distributions. North Holland Mathematical Studies 84, 1984.
[24] Robert Dautray and Jacques-Louis Lions. Analyse mathmatique et calcul
numrique pour les sciences et les techniques. Collection du CEA. Masson,
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[43] Robert E. Greene and Steven G. Krantz. Function theory of one complex
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[46] Godfrey H. Hardy. A mathematicians apology. Cambridge University
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[48] Bruce J. Hunt. Rigorous discipline: Oliver Heaviside versus the mathematicians. In Peter Dear, editor, The literary structure of scientific argument,
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Portraits
Airy Sir George . . . . . . . . . . . .
Alembert Jean le Rond d . . . . . .
Bayes Thomas . . . . . . . . . . . . .
Bienaym Jules . . . . . . . . . . . .
Bolzano Bernhard . . . . . . . . . .
Borel mile . . . . . . . . . . . . . .
Cauchy Augustin . . . . . . . . . .
Christoffel Elwin . . . . . . . . . .
Dirac Paul . . . . . . . . . . . . . .
Dirichlet Gustav . . . . . . . . . .
Euler Leonhard . . . . . . . . . . .
Fourier Joseph . . . . . . . . . . . .
Fubini Guido . . . . . . . . . . . . .
Gauss Carl . . . . . . . . . . . . . .
Gibbs Josiah . . . . . . . . . . . . . .
Heaviside Oliver . . . . . . . . . . .
Hilbert David . . . . . . . . . . . .
Jordan Camille . . . . . . . . . . . .
Kolmogorov Andrei Nikolaievitch
Kronecker Leopold . . . . . . . . .
Landau Edmund . . . . . . . . . . .
Laplace, marquis de . . . . . . . . .
Lebesgue Henri . . . . . . . . . . . .
Lvy Paul . . . . . . . . . . . . . . .
Liouville Joseph . . . . . . . . . . .
Minkowski Hermann . . . . . . . .
Poincar Henri . . . . . . . . . . .
Poisson Denis . . . . . . . . . . . .
Riemann Bernhard . . . . . . . . . .
Schmidt Erhard . . . . . . . . . . . .
Schwartz Laurent . . . . . . . . . .
Schwarz Hermann . . . . . . . . . .
Stokes George . . . . . . . . . . . .
Taylor Brook . . . . . . . . . . . . .
Tchebychev Pafnouti . . . . . . . .
Weierstrass Karl . . . . . . . . . . .
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42
415
518
547
581
58
88
163
186
171
39
268
80
557
311
193
257
119
515
445
580
333
61
558
104
448
475
172
138
258
222
161
472
31
548
582
Sidebars
Lebesgue measure for Borel sets . . . . . . . . . . .

A nonmeasurable set . . . . . . . . . . . . . . . . .
Differentiability of a function on R2 . . . . . . .
Pointwise convergence and Gibbs phenomenon .
Extension of a continuous linear operator . . . .
Integration of differential forms . . . . . . . . . .
Double connectedness of SO(3) and magic tricks
Riemann-Stieltjes integral . . . . . . . . . . . . . .
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60
70
134
276
294
488
506
552
Index
The boldface page numbers refer to the
main definitions, the italic ones to exercises. The family names in Small Capitals
refer to short biographies.
Symbols
(Hodge operator) . . . . . . . . . . . . . . . . . 482
(adjoint) . . . . . . . . . . . . . . . . . . . . . . . . . 378
(dual) . . . . . . . . . . . . . . . . . . . . . . . . 436, 463
(adjoint) . . . . . . . . . . . . . . . . . . . . . 389, 497
(convolution product) . . . . . . . . . . . . . 214
c.d.
. . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 554
c.p.
. . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 554
c.a.s.
. . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 554
cv.s.
. . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 19
cv.u.
. . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 19
(rectangle function) . . . . . . . . . . 213, 279
. . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 333
X . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 186
. . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 414
A (closure) . . . . . . . . . . . . . . . . . . . . . . . . 391
| . . . . . . . . . . . . . . . . . . . . . . . . . . . 260, 380
| . . . . . . . . . . . . . . . . . . . . . . . . . . . 260, 379
x| . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 382
|x . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 384
p| . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 383
|p . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 383
-algebra . . . . . . . . . . . . . . . . . . . . . . . . . . . 57
generated . . . . . . . . . . . . . . . . . . 58, 513
by an r.v. . . . . . . . . . . . . . . . . . . . 537
independent s . . . . . . . . . . . . 519, 537
(un vn ) . . . . . . . . . . . . . . . . . . . . . 580
. . . . . . . . . . . . . . . . . . . . . . . . . . . . . 210, 439
(exterior product) . . . . . . . . 464, 467, 468
Laurent Pierre . . . . . . . . . . . . . . . . . . . . . 112
A
a.e. . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 61
Abelian group . . . . . . . . . . . . . . . . . . . . . 489
Absolute convergence . . . . . . . . . . . . . . . . 29
Accumulation point . . . . . . . . . . . . . . . . 106
Adjoint . . . . . . . . . . . . . . . . . . . . . . . 378, 389
Airy George (sir) . . . . . . . . . . . . . . . . . . . . 42
Airy integral . . . . . . . . . . . . . . . . . . . . . . . . 42
dAlembert Jean le Rond . . . . . . . . . . . 415
dAlembertian . . . . . . . . . . . . . . . . . . . . . . 414
Algebra
-algebra . . . . . . . . . . . . . . . . . . . . . . 57
convolution . . . . . . . . . . . . . . . . 236
exterior . . . . . . . . . . . . 469, 463471
Lie . . . . . . . . . . . . . . . . . . . . . . . . . 496
Almost
complete system . . . . . . . . . . . . . . . 518
everywhere . . . . . . . . . . . . . . . . . . . . . 61
surely . . . . . . . . . . . . . . . . . . . . . . . . 515
Analytic
continuation . . . . . . . . . . . . . . 144, 366
function . . . . . . . . . . . . . . . . . . . 35, 101
signal . . . . . . . . . . . . . . . . . . . . 365, 366
Antihermitian function . . . . . . . . . . . . . 285
Antiholomorphic function . . . . . . . 92, 157
Argument . . . . . . . . . . . . . . . . . . . . . . . . . . 93
Asymptotic expansion . . . . . . . . . . . . . . . . 37
Autocorrelation . . . . . . . . . . . . . . . . 368, 370
Average power of a function . . . . . . . . . 370
Axiom of choice . . . . . . . . . . . . . . . . . . . . 70
B
B(R) . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 58
B(n, p) . . . . . . . . . . . . . . . . . . . . . . . . . . . 524
Ball (open/closed ) . . . . . . . . . . . . . . . . 578
Banach fixed point theorem . . . . . . . . . . 15
Basis
algebraic . . . . . . . . . . . . . . . . . . . 250
dual . . . . . . . . . . . . . . . . . . . 436, 463
generalized . . . . . . . . . . . . . . . . . 385
Hilbert . . . . . . . . . . . . . . . . . . . . . 258
Bayes Thomas . . . . . . . . . . . . . . . . . . . . . 518
Bayes formula . . . . . . . . . . . . . . . . . . 517, 518
Bernoulli Jacques . . . . . . . . . . . . . . . . . 555
Bernoulli distribution . . . . . . . . . . . . . . . 524
Bertrand function . . . . . . . . . . . . . . . . . . . 74
Bessel
function . . . . . . . . . . . . . . . . . . 320, 324
inequality . . . . . . . . . . . . . . . . . . . . 256
Bienaym Jules . . . . . . . . . . . . . . . . . . . . . 547
Bienaym identity . . . . . . . . . . . . . . . . . . 545
Bienaym-Tchebychev inequality . . . . . . 547
Bilateral Laplace transform . . . . . . . . . . 332
Binomial
distribution . . . . . . . . . . . . . . . . . . . 524
law . . . . . . . . . . . . . . . . . . . . . . . . . . . 616
Bolzano Bernhard . . . . . . . . . . . . . . . . . 581
Borel mile . . . . . . . . . . . . . . . . . . . . . . . . 58
Borel set . . . . . . . . . . . . . . . . . . . . . . . . . . . . 58
Bounded (continous) operator . . . . . . . 390
632
Bra . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 380
| . . . . . . . . . . . . . . . . . . . . . . 260, 380
x| . . . . . . . . . . . . . . . . . . . . . . . . . . . 382
p| . . . . . . . . . . . . . . . . . . . . . . . . . . 383
generalized . . . . . . . . . . . . . . . . . 383
Branch point . . . . . . . . . . . . . . . 111, 136, 139
Bromwich contour . . . . . . . . . . . . . . . . . 337
Brownian motion . . . . . . . . . . . . . . 170, 425
Buffon George Leclerc, comte de . . . . 549
C
cv.s.
.
.
.
.
.
.
.
.
.
.
.
.
.
.
.
. . . . . . . . . . . . . . . . . 19
cv.u.
. . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 19
C . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 146
C ([a, b ]) . . . . . . . . . . . . . . . . . . . . . . . . . . 577
Casorati-Weierstrass theorem . . . . . . . . . 109
Cauchy
criterion . . . . . . . . . . . . . . . . . . . . 13, 23
distribution . . . . . . . . . . . . . . . 529, 562
formula . . . . . . . . . . . . . . . . 45, 99, 101
law . . . . . . . . . . . . . . . . . . . . . . . . . . . 616
-Lipschitz theorem . . . . . . . . . . . . 46
principal value pv(1/x)188, 189, 224
Fourier transform of the . . . 304
problem . . . . . . . . . . . . . . . . . . 238, 346
product . . . . . . . . . . . . . . . . . . . . . . 215
-Riemann equations . . . . 88, 90, 92
-Schwarz inequality . . . . . . . . . . 252
sequence . . . . . . . . . . . . . . . . . . . . . . . 13
theorem . . . . . . . . . . 93, 97, 98, 98, 598
Cauchy Augustin-Louis . . . . . . . . . . . . . . 88
Causal
function . . . . . . . . . . . . . . . . . . 331, 366
system . . . . . . . . . . . . . . . . . . . . . . . . 219
Causality . . . . . . . . . . . . . . . . . . . . . . . . . . 550
Cavendish experiment . . . . . . . . . . . . . . 484
Centered random variable . . . . . . . . . . . 532
Central limit theorem . . . . . . . . . . . . . . . 558
Characteristic function A . . . . . . . . . . . 63
Characteristic function X . . . . . . . . . . 541
Charge (density of ) . . . . . . . . . . . . . . . 199
Choice (Axiom of ) . . . . . . . . . . . . . . . . 70
Christoffel Elwin . . . . . . . . . . . . . . . . . 163
Circular lens . . . . . . . . . . . . . . . . . . . . . . . 320
Class S . . . . . . . . . . . . . . . . . . . . . . . . . . . 397
Class (complete of events) . . . . . . . . . 513
Classical limit . . . . . . . . . . . . . . . . . . . . . . . . 7
Clever trick . . . . . . . . . . . . . . . . . . . . . . . . 142
Closable operator . . . . . . . . . . . . . . . . . . . 391
Closed . . . . . . . . . . . . . . . . . . . . . . . . . . . . 578
differential form . . . . . . . . . . . . . . 474
operator . . . . . . . . . . . . . . . . . . . . . . 390
Closure . . . . . . . . . . . . . . . . . . . . . . . . . . . . 579
of an operator . . . . . . . . . . . . . . . . 391
Index
of a subset . . . . . . . . . . . . . . . . . . . . 579
Coefficients (Fourier ) . . . . . . . . . 258, 264
Coherence
function . . . . . . . . . . . . . . . . . . . . . . 369
spatial . . . . . . . . . . . . . . . . . . . . . . . 43
temporal . . . . . . . . . . . . . . . . . . . 371
Comb . . . . . . . . . . . . . . . . . . see Dirac, comb
Commutative group . . . . . . . . . . . . . . . . 489
Commutatively convergent . . . . . . . . . . . 26
Compact . . . . . . . . . . . . . . . . . . . . . . . . . . 576
Complement . . . . . . . . . . . . . . . . . . . . . . . 512
Complete
measure space . . . . . . . . . . . . . . . . 61
class of events . . . . . . . . . . . . . . . . . 513
normed vector space . . . . . . . . . . . . 14
Complex logarithm . . . . . . . . . . . . . 135, 136
Complex velocity (of a fluid) . . . . . . . . 167
Conditional probability . . . . . . . . . . . . . 517
Conditionally convergent series . . . . . . . 26
Conformal map . . . . . . . . . . . . . . . . . . . . 156
Conformally equivalent . . . . . . . . . . . . . 157
Connected
doubly . . . . . . . . . . . . . . . . . . . . . 501
simply . . . . . . . . . . . . . . . . . . 114, 575
space . . . . . . . . . . . . . . . . . . . . . . . . . 574
Constants (structure ) . . . . . . . . . . . . . 496
Continuation (analytic ) . . . . . . . 144, 366
Continuity
of convolution . . . . . . . . . . . . . . . . 235
of differentiation . . . . . . . . . . . . . . 230
of a linear functional . . . . . . . . . . 184
of a map . R. . . . . . . . . . . . . . . . . . . . 574
under the sign . . . . . . . . . . . . . . . 77
Contour
Bromwich . . . . . . . . . . . . . . . . . . 337
integration . . . . . . . . . . . . . . . . . . . 94
Contractible open set . . . . . . . . . . . . . . . 474
Contraction . . . . . . . . . . . . . . . . . . . . . . . . . 15
Contraction of indices . . . . . . . . . . . . . . 454
Contravariant coordinates . . . . . . . . . . . 434
Convergence
absolute . . . . . . . . . . . . . . . . . . 23, 29
almost everywhere . . . . . . . . . . . 554
almost sure . . . . . . . . . . . . . . . . . 554
commutative . . . . . . . . . . . . . . . . . 26
conditional . . . . . . . . . . . . . . . . . . 26
in D . . . . . . . . . . . . . . . . . . . . . . . . . 183
in distribution . . . . . . . . . . . . . . . . 554
dominated theorem . . . . . . . . . . . 75
in measure . . . . . . . . . . . . . . . . . . . . 554
normal . . . . . . . . . . . . . . . . . . . . . . 23
pointwise
of a Fourier series . . . . . . . . . . . 268
of a sequence of functions . . . . 19
Index
of a series of functions . . . . . . . 29
in probability . . . . . . . . . . . . . . . . . 554
radius of . . . . . . . . . . . . . . . . . . . . 34
in S . . . . . . . . . . . . . . . . . . . . . . . . 290
of a sequence . . . . . . . . . . . . . . . . . . 12
of a series . . . . . . . . . . . . . . . . . . . . . 23
simple . . . . . . . . . . . see pointwise
to . . . . . . . . . . . . . . . . . . . . . . . . 13
uniform
of a double sequence . . . . . . . . . 16
of a Fourier series . . . . . . . . . . . 269
of a sequence of functions . . . . 19
of a series of functions . . . . . . . 29
weak . . . . . . . . . . . . . . . . . . . . . . . 230
Convexity . . . . . . . . . . . . . . . . . . . . . . . . . 575
Convolution
algebra . . . . . . . . . . . . . . . . . . . . . . . 236
of causal functions . . . . . . . . . . . . 339
continuity of . . . . . . . . . . . . . . . 235
discrete . . . . . . . . . . . . . . . . . . . . . 220
of distributions . . . . . . . . . . . . . . . . 214
and Fourier transform . . . . . . . . . 292
of functions . . . . . . . . . . . . . . . 211, 270
inverse . . . . . . . . . . . . . . . . . . . . . . . 236
regularization by . . . . . . . . . . . . 235
Coordinate forms . . . . . . . . . . . . . . . . . . 436
Coordinates
change . . . . . . . . . . . . . . . . . . . . . . . 455
contravariant . . . . . . . . . . . 434, 451
covariant . . . . . . . . . . . . . . . 437, 449
curvilinear . . . . . . . . . . . . . . . . . . 455
Correlation . . . . . . . . . . . . . . . . . . . . . . . . 535
Correlation function . . . . . . . . . . . . . . . . 368
Coulomb potential . . . . . . . . . . . . . . . . . 166
Fourier transform of the . . . . . 306
Laplacian of the . . . . . . . . . . . . . 208
Covariance . . . . . . . . . . . . . . . . . . . . . . . . 535
Covariant coordinates . . . . . . . . . . . . . . . 437
Criterion
Cauchy . . . . . . . . . . . . . . . . . . . 13, 23
Rayleigh . . . . . . . . . . . . . . . . . . . . 321
Current (density of ) . . . . . . . . . . . . . . 199
Curve . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 93
Curvilinear coordinates . . . . . . . . . . . . . 455
Cut . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 136
Cyclic group . . . . . . . . . . . see Permutations
D
D . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 182
is dense in D . . . . . . . . . . . . . . . . 235
D+ , D . . . . . . . . . . . . . . . . . . . . . . . . . . . 236
. . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 196
D . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 183
d (exterior derivative) . . . . . . . . . . . . . . . 470
633
. . . . . . . . . . . . . . . . . . . . . . . . . . . . . 185, 186
dx . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 91, 134
dx i dx j . . . . . . . . . . . . . . . . . . . . . . . . . . 464
dz, dz . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 91
/ z, / z . . . . . . . . . . . . . . . . . . . . . . . . . 92
dAlembertian . . . . . . . . . . . . . . . . . . . . . . 414
dAlembert Jean le Rond . . . . . . . . . . . 415
Debye Petrus . . . . . . . . . . . . . . . . . . . . . . 148
Debye screen, potential . . . . . . . . . . . . . . 325
Deceit . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 571
Degree of a representation . . . . . . . . . . . 492
Dense subset . . . . . . . . . . . . . . . . . . . 574, 579
Density
of charge and current . . . . . . . . . . 199
of D in D . . . . . . . . . . . . . . . . . . . 235
probability . . . . . . . . . . . . . . . . . 526
joint . . . . . . . . . . . . . . . . 534, 539
of S in L2 . . . . . . . . . . . . . . . . . . . 290
spectral . . . . . . . . . . . . . . . . . 368, 370
of a subset . . . . . . . . . . . . . . . . 574, 579
Derivation (continuity of ) . . . . . . . . . 230
Derivative
of a discontinuous function . . . . 201
of a distribution . . . . . . . . . . . . . . . 192
exterior . . . . . . . . . . . . . . . . . . . . 470
C 1 -diffeomorphism . . . . . . R. . . . . . . . . . 455
Differentiability under the sign . . . . . 78
Differential
form . . . . . . . . . . . . . . . . . 469, 463483
integral of a . . . . . . . . . . . . . . 471
of a function . . . . . . . . . 134, 586, 587
Diffraction . . . . . . . . . . . . . . . . . . . . . . . . . 314
Diffuse
mesure . . . . . . . . . . . . . . . . . . . . . . . . 60
probability . . . . . . . . . . . . . . . . . . . 525
Dilation of a distribution . . . . . . . . . . . 190
Dini Ulisse . . . . . . . . . . . . . . . . . . . . . . . . . 22
Dinis theorems . . . . . . . . . . . . . . . . . . . . . 22
Dirac
comb X . . . . . . . . . . . . . . . . . 186, 307
Fourier transform of . . . . . . 307
curvilinear distribution . . . . . . 195
distribution . . . . . . . . . 185, 186, 194
Laplace transform of . . . . . . 342
distribution . . . . . . . . . . . . . . . . 196
measure . . . . . . . . . . . . . . . . . . 514, 552
sequence . . . . . . . . . . . . . . . . . . . . . . 231
surface distribution . . . . . . . . . . 195
Dirac Paul . . . . . . . . . . . . . . . . . . . . . . . . 186
Direct product . . . . . . . . . . . . . . . . . 209, 210
Dirichlet
function . . . . . . . . . . . . . . . . . . . 62, 279
problem . . . . . . . . . . . . . . . . . . . . . . 170
634
for a disc . . . . . . . . . . . . . . . . . . 172
on a half-plane . . . . . . . . . . . . . 175
for a strip . . . . . . . . . . . . . . . . . . 175
theorem . . . . . . . . . . . . . . . . . . . . . . 268
Dirichlet Gustav Lejeune . . . . . . . . . . . 171
Dispersion relation . . . . . . . . . . . . . . . . . 227
Distance . . . . . . . . . . . . . . . . . . . . . . . . . . . 578
Distribution . . . . . . . . . . . . . . . . . . . . . . . 183
Bernoulli . . . . . . . . . . . . . . . . . . . 524
binomial . . . . . . . . . . . . . . . . . . . 524
Cauchy . . . . . . . . . . . . . . . . 529, 562
convergence in . . . . . . . . . . . . . . 554
curvilinear Dirac . . . . . . . . . . . . 195
(Dirac derivative) . . . . . . . . . . . 196
derivative of a . . . . . . . . . . . . . . . 192
dilation of a . . . . . . . . . . . . . . . . 190
Dirac . . . . . . . . . . . . 185, 186, 194
Laplace transform of . . . . . . 342
function . . . . . . . . . . . . . . . . . . . . . . 524
joint . . . . . . . . . . . . . . . . . . . . 534
Gaussian . . . . . . . . . . . . . . . . . . . 557
Heaviside . . . . . . . . . . . . . . . . . . . 193
marginal . . . . . . . . . . . . . . . . . . . 534
normal . . . . . . . . . . . . . . . . . 515, 557
Poisson . . . . . . . . . . . . 530, 530, 547
probability . . . . . . . . . . . . . . . . . 522
regular . . . . . . . . . . . . . . . . . . . . . 184
regularization of a . . . . . . . . . . . 234
sgn (sign) . . . . . . . . . . . . . . . . . 203
singular . . . . . . . . . . . . . . . . . . . . 185
support of a . . . . . . . . . . . . . . . . 187
surface Dirac . . . . . . . . . . . . . . . . 195
tempered . . . . . . . . . . . . . . . . . . . 300
translate of a . . . . . . . . . . . . . . . . 189
transpose of a . . . . . . . . . . . . . . . 190
Domain . . . . . . . . . . . . . . . . . . . . . . . . . . . . 90
Dominated
convergence theorem . . . . . . . . . . . 75
sequence by another . . . . . . . . . 580
Doubly connected . . . . . . . . . . . . . . . . . . 501
Dual
basis . . . . . . . . . . . . . . . . . . . . . 436, 463
Hodge . . . . . . . . . . . . . . . . . . . . . 482
of a vector space . . . . . . . . . . 436, 463
Duality (metric ) . . . . . . . . . . . . . . . . . . 449
E
E . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 236
() . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 465
. . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 483
E . . . . . . . . . . . . . . . . . . . . . . . . . . . . 436, 463
Egorov theorem . . . . . . . . . . . . . . . . . . . . . 69
Eigenvalue . . . . . . . . . . . . . . . . . . . . . 391, 403
Index
generalized . . . . . . . . . . . . . 392, 403
Eigenvector . . . . . . . . . . . . . . . . . . . . . . . . 391
Einstein Albert . . . . . . . . . . . . . . . . . . . . 435
Einstein convention . . . . . . . . . . . . . . . . 435
Electromagnetic field
dynamics of the without sources
348
Green function of the . . . . . . . . 414
Electromagnetism . . 228, 325, 348, 414422
Electrostatics . . . . . . . . . . . . . . . . . . . 165, 194
Element (matrix ) . . . . . . . . . . . . . . . . . 398
Endomorphism
hermitian . . . . . . . . . . . . . . . . . . 379
normal . . . . . . . . . . . . . . . . . . . . . 379
self-adjoint . . . . . . . . . . . . . . . . . 379
symmetric . . . . . . . . . . . . . . . . . . 379
Energy (of a signal) . . . . . . . . . . . . . . . . . 368
Entire function . . . . . . . . . . . . . . . . . . . . 103
Equation
heat . . . . . . . . . . . . . . . . . . . 240, 422
Klein-Gordon . . . . . . . . . . . . . . . 429
Maxwell . . . . . . . . . . . . . . . . 481, 482
Poisson . . . . . . . . . . . . . . . . . 164, 217
Schrdinger . . . . . . . . . . . . . . . . . 427
Equivalent
conformally open sets . . . . . . . . 157
norms . . . . . . . . . . . . . . . . . . . . . . . . 581
paths . . . . . . . . . . . . . . . . . . . . . . . . . . 94
sequences . . . . . . . . . . . . . . . . . . . . . 580
Essential singularity . . . . . . . . . . . . . . . . . 110
Euler Leonhard . . . . . . . . . . . . . . . . . . . . 39
Euler function . . . . . . . . . . . . . . . . . . . . . 154
Event . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 512
incompatible s . . . . . . . . . . . . . . . 513
realized . . . . . . . . . . . . . . . . . . . . . 513
Exact differential form . . . . . . . . . . . . . . 474
Expansion
asymptotic series . . . . . . . . . . . . . 37
Fourier series . . . . . . . . . . . . . . . 265
perturbative . . . . . . . . . . . . . . . . . . 38
power series . . . . . . . . . . . . . . . 34, 35
Taylor . . . . . . . . . . . . . . . . . . . . . . . 36
Expectation . . . . . . . . . . . . . . . . . . . . 527, 528
Extension
of a continuous operator . . . . . . . 294
of an operator on H . . . . . . . . . . . 388
Exterior
1-form . . . . . . . . . . . . . . . . . . . . . . . 463
2-form . . . . . . . . . . . . . . . . . . . . . . . 464
k-form . . . . . . . . . . . . . . . . . . . . . . . 465
algebra . . . . . . . . . . . . . . . 469, 463471
derivative . . . . . . . . . . . . . . . . . . . . . 470
product . . . . . . . . . . . . . . . . . . 467, 468
Index
F
F [ f ] . . . . . . . . . . . . . . . . . . . . . . . . . . . . 278
Faithful representation . . . . . . . . . . . . . . 492
Faltung theorem . . . . . . . . . . . . . . . . . . . . . 292
Family (free/generating ) . . . . . . . . . . . 250
Faraday tensor . . . . . . . . . . . . . . . . . . . . . 481
Fejr sums . . . . . . . . . . . . . . . . . . . . . 270, 313
Feynman Richard . . . . . . . . . . . . . . . . . . 432
Feynman propagator . . . . . . . . . . . . . . . . 432
Field
electromagnetic
dynamics of the without sources
348
Green function of the . . . . . 414
transverse/longitudinal . . . . . . 358
Finite part fp(1/x k ) . . . . . . . . . 241, 324, 345
Finite power functions . . . . . . . . . . . . . . 370
Fischer (Riesz- theorem) . . . . . . . . . . . . 67
Fixed point . . . . . . . . . . . . . . . . . . . . . . . . . 15
Form
coordinate . . . . . . . . . . . . . . . . . . 436
differential . . . . . . . . . 469, 463483
closed . . . . . . . . . . . . . . . . . . . 474
exact . . . . . . . . . . . . . . . . . . . . 474
integral of a . . . . . . . . . . . . . . 471
exterior 1-form . . . . . . . . . . . . . . . . 463
exterior 2-form . . . . . . . . . . . . . . . . 464
exterior k-form . . . . . . . . . . . . . . . . 465
linear . . . . . . . . . . . . . . . . . . 436, 463
volume . . . . . . . . . . . . . . . . . . . . . 467
Formula
Bayes . . . . . . . . . . . . . . . . . . . 517, 518
Cauchy . . . . . . . . . . . . . . . 45, 99, 101
Green . . . . . . . . . . . . . . . . . . . . . . 206
Green-Ostrogradski . . . . . . 205, 477
Gutzmer . . . . . . . . . . . . . . . . . . . . 45
Poincar . . . . . . . . . . . . . . . . 516, 606
Poisson summation . . . . . . 271, 309
Stirling . . . . . . . . . . . . . . . . . . . . . 154
Stokes . . . . . . . . . . . . . . . . . . . . . . 473
Taylor . . . . . . . . . . . . . . . . . . . . . . . 31
Taylor-Lagrange . . . . . . . . . . . . 31, 32
Taylor-Young . . . . . . . . . . . . . . . . . 32
Four-vector . . . . . . . . . . . . . . . . . . . . . . . . 200
Fourier
coefficients . . . . . . . . . . . . . . . 258, 264
partial series . . . . . . . . . . . . . . . . 258
series . . . . . . . . . . . . . . . . 258, 264, 265
transform
of X . . . . . . . . . . . . . . . . . . . . . . 307
computation by residues . . . . . 120
conjugate transform . . . . . . . 278
of . . . . . . . . . . . . . . . . . . . . . . . 303
of a distribution . . . . . . . . . . . . 300
635
of a function . . . . . . . . . . . . . . . 278
of the gaussian . . . . . . . . . 126, 295
of H . . . . . . . . . . . . . . . . . . . . . . 304
inverse of the 282, 284, 291, 306
of the laplacian . . . . . . . . . . . . . 305
of the lorentzian 1/(1 + t 2 ) . . . 84
of pv(1/x) . . . . . . . . . . . . . . . . . 304
in Rn . . . . . . . . . . . . . . . . . . . . . . 358
sine and cosine . . . . . . . . . . . 295
Fourier Joseph . . . . . . . . . . . . . . . . . . . . 268
fp(1/x k ) . . . . . . . . . . . . . . . . . . . . . . . 324, 345
Fraunhofer approximation . . . . . . . . . . . 314
Free family . . . . . . . . . . . . . . . . . . . . . . . . 250
Frequencies (Matsubara ) . . . . . . . . . . . 128
Fubini Guido . . . . . . . . . . . . . . . . . . . . . . . 80
Fubini-Lebesgue theorem . . . . . . . . . . . . . 79
Fubini-Tonelli theorem . . . . . . . . . . . . . . . 80
Function
analytic . . . . . . . . . . . . . . . . . . 35, 101
antihermitian . . . . . . . . . . . . . . . 285
antiholomorphic . . . . . . . . . 92, 157
autocorrelation . . . . . . . . . . 368, 370
Bertrand (t log t) . . . . . . . . . . . 74
Bessel . . . . . . . . . . . . . . . . . . 320, 324
causal . . . . . . . . . . . . . . . . . . 331, 366
characteristic . . . . . . . . . . . . . 63, 541
coherence . . . . . . . . . . . . . . . . . . 369
Dirac function . . . . . . . . . . 185, 186
Dirichlet . . . . . . . . . . . . . . . . 62, 279
distribution . . . . . . . . . . . . . . . . . 524
joint . . . . . . . . . . . . . . . . 534, 538
entire . . . . . . . . . . . . . . . . . . . . . . 103
Euler . . . . . . . . . . . . . . . . . . . . . . . 154
finite power . . . . . . . . . . . . . . . . . 370
Green . . . . . . 165, 236, 408, 407432
of the dAlembertian . . . . 414, 417
of the harmonic oscillator . . . 409
of the heat equation . . . . 423, 424
harmonic . . . . . . . . . . . 139, 139144
Heaviside . . . . . . . . . . . . . . . . 62, 193
hermitian . . . . . . . . . . . . . . . . . . 285
holomorphic . . . . . . . . . . . . . . . . 90
integrable . . . . . . . . . . . . . . . . . . . . 64
intercorrelation . . . . . . . . . . . . . 369
locally integrable . . . . . . . . . . . . 184
measurable . . . . . . . . . . . . . . . . . . 62
meromorphic . . . . . . . . . . . . . . . 110
multivalued . . . . . . . . . . . . . 135139
rapidly decaying . . . . . . . . . . . . . 288
rectangle (x) . . . . . . . . . . 213, 279
self-coherence . . . . . . . . . . . . . . . 369
simple . . . . . . . . . . . . . . . . . . . . . . 63
slowly increasing . . . . . . . . . . . . 301
test . . . . . . . . . . . . . . . . . . . . . . . . 182
636
transfer . . . . . . . . . . . . . . . . . . . . 356
Functional . . . . . . . . . . . . . . . . . . . . . . . . . 182
continuity of a linear . . . . . . . . 184
G
g . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 447
b
g . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 451
g . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 451
g . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 447
(z) (Euler function) . . . . . . . . . . . . . . . 154
Gauge . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 590
Gauss Carl . . . . . . . . . . . . . . . . . . . . . . . . 557
Gaussian . . . . . . . . . . . . . . . . . . . . . . . . . . 279
distribution . . . . . . . . . . . . . . . 515, 557
Fourier transform . . . . . . . . . 126, 295
Gelfand triple . . . . . . . . . . . . . . . . . . . . . . 383
Generated (-algebra ) . . . . . . . . . . 58, 513
Generators (infinitesimal ) . . . . . . . . . 495
Gibbs Josiah . . . . . . . . . . . . . . . . . . . . . . . 311
Gibbs phenomenon . . . . . . . . . . . . . . . . . 311
Graph of a linear operator . . . . . . . . . . 390
Green
-Ostrogradski formula . . . 205, 477
formula . . . . . . . . . . . . . . . . . . . . . . 206
function . . . . . 165, 236, 408, 407432
of the dAlembertian . . . . 414, 417
of the harmonic oscillator . . . 409
of the heat equation . . . . 423, 424
theorem . . . . . . . . . . . . . . . . . . . . . . 106
Green George . . . . . . . . . . . . . . . . . . . . . 105
Green-Riemann theorem . . . . . . . . . . . . 105
Group . . . . . . . . . . . . . . . . . . . . . . . . . . . . 489
abelian (commutative) . . . . . . . 489
linear representation . . . . . . . . . . . 492
one-parameter . . . . . . . . . . . . . . . 495
of rotations . . . . . . . . . . . . . . . . . . . 492
Gutzmer formula . . . . . . . . . . . . . . . . . . . . 45
H
H (x) . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 193
Half-life . . . . . . . . . . . . . . . . . . . . . . . . . . . 533
Hankel transform . . . . . . . . . . . . . . . . . . 324
Harmonic function . . . . . . . . . 139, 139144
Harmonic oscillator . . . . . . . . . . . . 239, 409
Heat
equation . . . . . . . . . . . . . . . . . 240, 422
kernel . . . . . . . . . . . . . . . . 240, 423, 424
operator . . . . . . . . . . . . . . . . . . . . . . 240
Heaviside
distribution . . . . . . . . . . . . . . . . . . . 193
Fourier transform of H . . . . . . . . 304
function . . . . . . . . . . . . . . . . . . . 62, 193
Laplace transform of H . . . . . . . . 334
Heaviside Oliver . . . . . . . . . . . . . . . . . . . 193
Index
Heisenberg uncertainty relations . 363, 404
Hermite polynomials . . . . . . . . . . . . . . . 263
Hermitian
endomorphism . . . . . . . . . . . . . . . . 379
function . . . . . . . . . . . . . . . . . . . . . . 285
operator . . . . . . . . . . . . . . . . . . . . . . 394
product . . . . . . . . . . . . . . . . . . . . . . 251
Hilbert
basis . . . . . . . . . . . . . . . . . . . . . . . . . 258
infinite hotels . . . . . . . . . . . . . . 12, 330
space . . . . . . . . . . . . . . . . . . . . . . . . . 257
transform . . . . . . . . . . . . . . . . . 227, 366
Hilbert David . . . . . . . . . . . . . . . . . . . . . 257
Hodge operator . . . . . . . . . . . . . . . . . . . . 482
Holomorphic function . . . . . . . . . . . . . . . 90
Holomorphically simply connected . . . 114
Homeomorphic open sets . . . . . . . . . . . 157
Homeomorphism . . . . . . . . . . . . . . . . . . 157
Homotopy . . . . . . . . . . . . . . . . . . . . . . . . . 575
I
I.r.v. . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 537
Identity
Bienaym . . . . . . . . . . . . . . . . . . . 545
parallelogram . . . . . . . . . . . . . . . 253
Parseval . . . . . . . . . . . . . . . . 259, 265
Parseval-Plancherel . . . . . . . . . . . 291
Poisson summation formula . . . . 271
Taylor-Lagrange . . . . . . . . . . . . . . 31
Image . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 332
Impulse response . . . . . . . . . . . . . . . . . . . 218
Independent
-algebras . . . . . . . . . . . . . . . . . . . . . 519
events . . . . . . . . . . . . . . . . . . . . . . . . 519
random variables . . . . . . . . . . . . . . 537
Indices
contracted . . . . . . . . . . . . . . . . . . 454
contravariant . . . . . . . . . . . 434, 451
covariant . . . . . . . . . . . . . . . 437, 449
Inequality
Bessel . . . . . . . . . . . . . . . . . . . . . . 256
Bienaym-Tchebychev . . . . . . . . 547
Cauchy-Schwarz . . . . . . . . . . . . . 252
Heisenberg . . . . . . . . . . . . . 363, 404
Minkowski . . . . . . . . . . . . . . . . . 253
Taylor-Lagrange . . . . . . . . . . . . . . 32
Infinite hotels . . . . . . . . . . . . . . . . . . . 12, 330
Infinitesimal generators . . . . . . . . . . . . . 495
Integrable function . . . . . . . . . . . . . . . 64, 64
Integral
of a differential form . . . . . . . . . . 471
Lebesgue . . . . . . . . . . . . . . . . . . . . 64
Riemann . . . . . . . . . . . . . . . . . 51, 53
Riemann-Stieltjes . . . . . . . . 529, 552
Index
Intercoherence function . . . . . . . . . . . . . 369
Intercorrelation . . . . . . . . . . . . . . . . . . . . 369
Interior of a subset . . . . . . . . . . . . . . . . . 578
Inverse (convolution ) . . . . . . . . . . . . . 236
Inversion
of the Fourier transform282, 284, 306
of the Laplace transform . . . . . . . 338
J
J0 (x), J1 (x) . . . . . . . . . . . . . . . . . . . . . . . . 320
Jacobian . . . . . . . . . . . . . . . . . . . . . . . . 81, 588
Joint
distribution function . . . . . . . . . . 534
probability density . . . . . . . . 534, 539
Jordan Camille . . . . . . . . . . . . . . . . . . . . 119
Jordans lemmas . . . . . . . . . . . . . . . . 117, 118
Joukovski Nicola . . . . . . . . . . . . . . . . . . 158
Joukovski mapping . . . . . . . . . . . . . . . . . 158
K
Kernel
Dirichlet . . . . . . . . . . . . . . . . . . . 270
Fejr . . . . . . . . . . . . . . . . . . . . . . . 270
heat . . . . . . . . . . . . . . . 240, 423, 424
Poisson . . . . . . . . . . . . . . . . . . . . . 172
Schrdinger . . . . . . . . . . . . . . . . 428
Ket . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 379
| . . . . . . . . . . . . . . . . . . . . . . 260, 379
|x . . . . . . . . . . . . . . . . . . . . . . . . . . . 384
|p . . . . . . . . . . . . . . . . . . . . . . . . . . 383
Khintchine (Wiener- theorem) . . . . . 371
Kirchhoff integral . . . . . . . . . . . . . . 244, 349
Klein-Gordon equation . . . . . . . . . . . . . 429
Kolmogorov Andrei Nikolaievitch . . . 515
Kramers-Kronig relations . . . . . . . . . . . . 227
Kronecker Leopold . . . . . . . . . . . . . . . . 445
Kronecker symbol . . . . . . . . . . . . . . . . . . 445
L
L . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 280
2 (E ) . . . . . . . . . . . . . . . . . . . . . . . . . . . . 464
k (E ) . . . . . . . . . . . . . . . . . . . . . . . . . . . . 466
L2 (R) . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 263
L2 [0, a] . . . . . . . . . . . . . . . . . . . . . . . 262, 264
L1 . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 280
L1loc (Rn ) . . . . . . . . . . . . . . . . . . . . . . . . . . . 184
2 . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 261
Lagrange
multipliers . . . . . . . . . . . . . . . . . . . . 588
Taylor- formula . . . . . . . . . . . . . . 32
Lagrangian (Proca ) . . . . . . . . . . . . . . . . 484
Landau Edmund . . . . . . . . . . . . . . . . . . 580
Landau notation . . . . . . . . . . . . . . . . . . . 580
Laplace, Pierre Simon de . . . . . . . . . . . 333
637
Laplace transform . . . . . . . . . . 332, 331350
of distributions . . . . . . . . . . . . . . . 342
inversion of the . . . . . . . . . . . . . 338
Laplacian . . . . . . . . . . . . . . . . . . . . . . . . . . 139
Large numbers
strong law of . . . . . . . . . . . . . . . 556
weak law of . . . . . . . . . . . . . . . . . 555
Laurent series . . . . . . . . . . . . . . . . . . 111, 112
Law
binomial . . . . . . . . . . . . . . . . . . . 616
Cauchy . . . . . . . . . . . . . . . . . . . . . 616
Poisson . . . . . . . . . . . . . . . . . . . . . 616
strong of large numbers . . . . . . 556
weak of large numbers . . . . . . . 555
Lebesgue
integral . . . . . . . . . . . . . . . . . . . . . . . . 64
measure . . . . . . . . . . . . . . . . . . . . . . . 59
Riemann- lemma . . . . . . . . 267, 282
theorem . . . . . . . . . . . . . . . . . . . . . . . 75
Lebesgue Henri . . . . . . . . . . . . . . . . . . . . . 61
Lemma
Jordan . . . . . . . . . . . . . . . . . . . . 117, 118
Riemann-Lebesgue . . . . . . . . . 267, 282
Zorn . . . . . . . . . . . . . . . . . . . . . . . . . 249
Length of a path . . . . . . . . . . . . . . . . . . . . 95
Levi (Beppo- theorem) . . . . . . . . . . . . . 76
Levi-Civita tensor . . . . . . . . . . . . . . 483, 496
Lvy Paul . . . . . . . . . . . . . . . . . . . . . . . . . 558
Lie algebra . . . . . . . . . . . . . . . . . . . . . . . . . 496
Lifetime . . . . . . . . . . . . . . . . . . . . . . . . . . . 533
Lightning rod . . . . . . . . . . . . . . . . . . . . . . 169
Limit
classical . . . . . . . . . . . . . . . . . . . . . . 7
in D . . . . . . . . . . . . . . . . . . . . . . . . . 183
pointwise of a series . . . . . . . . . . 29
uniform of a series . . . . . . . . . . . 29
Linear representation . . . . . . . . . . . . . . . 492
Liouville Joseph . . . . . . . . . . . . . . . . . . . 104
Liouville theorem . . . . . . . . . . . . . . . 45, 103
Lipschitz (Cauchy- theorem) . . . . . . . . 46
Locally
finite . . . . . . . . . . . . . . . . . . . . . . . . . 110
integrable function . . . . . . . . . . . . 184
Longitudinal fields . . . . . . . . . . . . . . . . . 358
Lorentzian . . . . . . . . . . . . . . . . . . . . . 279, 284
M
M1,3 . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 448
Magic . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 506
Markov inequality . . . . . . . . . . . . . . . . . . 548
Matrix . . . . . . . . . . . . . . . . . . . . . . . . . . . . 593
element . . . . . . . . . . . . . . . . . . . . . . 398
jacobian . . . . . . . . . . . . . . . . . 81, 588
638
orthogonal . . . . . . . . . . . . . . . . . 493
Pauli . . . . . . . . . . . . . . . . . . . 498, 505
representation . . . . . . . . . . . . . . . . . 594
rotation . . . . . . . . . . . . . . . . . . . . 492
unitary . . . . . . . . . . . . . . . . . . . . . 497
Wronskian . . . . . . . . . . . . . . . . . . 409
Matsubara frequencies . . . . . . . . . . . . . . 128
Maximum modulus principle . . . . . . . . 104
Maximum principle . . . . . . . . . . . . . . . . 141
Maxwell equations . . . . . . . . . . . . . . 481, 482
Mean value
theorem . . . . . . . . . . . . . . . . . . . . . . 141
Mean value property . . . . . . . . . . . . . . . . 100
Measurable
function . . . . . . . . . . . . . . . . . . . . . . . 62
space . . . . . . . . . . . . . . . . . . . . . . . . . . 59
Measure . . . . . . . . . . . . . . . . . . . . . . . . . . . . 59
diffuse . . . . . . . . . . . . . . . . . . . . . . 60
Dirac . . . . . . . . . . . . . . . . . . . . . . 552
Lebesgue . . . . . . . . . . . . . . . . . 59, 60
Lebesgue exterior . . . . . . . . . . . . . 60
Median . . . . . . . . . . . . . . . . . . . . . . . . . . . 533
Meromorphic function . . . . . . . . . . . . . . 110
Metric . . . . . . . . . . . . . . . . . . . . . . . . . . . . 447
duality . . . . . . . . . . . . . . . . . . . . . . . 449
Minkowski . . . . . . . . . . . . . . . . . 448
Minkowski
inequality . . . . . . . . . . . . . . . . . . . . 253
pseudo-metric . . . . . . . . . . . . . . . . . 448
space . . . . . . . . . . . . . . . . . . . . . . . . . 448
Minkowski Hermann . . . . . . . . . . . . . . . 448
de Moivre Abraham . . . . . . . . . . . . . . . . 511
Moment . . . . . . . . . . . . . . . . . . . . . . . . . . . 532
of a couple of r.v. . . . . . . . . . . . . . 535
Momentum
average . . . . . . . . . . . . . . . . . . . . . 361
operator . 361, 388, 389, 391, 401, 402
representation . . . . . . . . . . . . . 360, 361
uncertainty . . . . . . . . . . . . . . . . . . . 361
Monty Hall paradox . . . . . . . . . . . . . . . . 560
Moreras theorem . . . . . . . . . . . . . . . . . . . 105
Multipliers (Lagrange ) . . . . . . . . . . . . 588
Multivalued function . . . . . . . . . . . . . . . 137
N
N (m, 2 ) . . . . . . . . . . . . . . . . . . . . . . . . . 557
Negligible
sequence compared to another 580
set . . . . . . . . . . . . . . . . . . . . . . . . 61, 515
Neighborhood . . . . . . . . . . . . . . . . . . . . . 573
Neumann problem . . . . . . . . . . . . . . . . . 167
Nonmonochromatic signal . . . . . . . . . . 372
Norm . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 577
equivalent s . . . . . . . . . . . . . . . . . 581
Index
hermitian . . . . . . . . . . . . . . . . . . 252
Normal distribution . . . . . . . . . . . . 515, 557
Normal endomorphism . . . . . . . . . . . . . 379
Normed vector space . . . . . . . . . . . . . . . . 577
O
O(3) . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 493
O(n ), o(n ) . . . . . . . . . . . . . . . . . . . . . . . 580
Observable . . . . . . . . . . . . . . . . . . . . . . . . 393
Open set . . . . . . . . . . . . . . . . . . . . . . 573, 578
conformally equivalent s . . . . . . 157
contractible . . . . . . . . . . . . . . . . . 474
star-shaped . . . . . . . . . . . . . . . . . . 475
Open/closed ball . . . . . . . . . . . . . . . . . . . 578
Operator
bounded (continuous) . . . . . . . 390
closable . . . . . . . . . . . . . . . . . . . . 391
closed . . . . . . . . . . . . . . . . . . . . . . 390
heat . . . . . . . . . . . . . . . . . . . . . . . . 422
hermitian (symmetric) . . . . . . . 394
Hodge . . . . . . . . . . . . . . . . . . . . . 482
momentum . . . . 361, 388, 389, 391,
401, 402
position . . . . . . . . 360, 391, 402, 403
Schrdinger . . . . . . . . . . . . . . . . . 427
self-adjoint . . . . . . . . . . . . . . . . . 394
Optics . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 314
Order of a pole . . . . . . . . . . . . . . . . . . . . 112
Original . . . . . . . . . . . . . . . . . . . . . . . . . . . 332
Orthogonal
matrix . . . . . . . . . . . . . . . . . . . . . . . 493
projection . . . . . . . . . . . . . . . . 255, 399
system . . . . . . . . . . . . . . . . . . . . . . . 253
Orthogonality . . . . . . . . . . . . . . . . . . . . . 253
Orthonormal system . . . . . . . . . . . . . . . . 253
Ostrogradski (Green- formula) . 205, 477
P
P (momentum operator) . . . . . . . . 388, 389
P . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 57, 512
(rectangle function) . . . . . . . . . . 213, 279
Parseval
-Plancherel identity . . . . . . . . . . 291
identity . . . . . . . . . . . . . . . . . . 259, 265
Partial Fourier series . . . . . . . . . . . . . . . . 258
Pascal Blaise . . . . . . . . . . . . . . . . . . . . . . 606
Path . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 93
integral on a . . . . . . . . . . . . . . . . . 94
length of a . . . . . . . . . . . . . . . . . . 95
winding number of a . . . . . . . . . 96
Pauli matrices . . . . . . . . . . . . . . . . . . . . . . 498
Percolation . . . . . . . . . . . . . . . . . . . . . . . . 169
Permutation . . . . . . . . . . . . . . . . . . . . 26, 465
Permutations . . . . . . . see Symmetric, group
Index
Ptanque . . . . . . . . . . . . . . . . . . . . . . . . . . . 181
fp(1/x k ) . . . . . . . . . . . . . . . . . . . . . . . . . . . 241
Phenomenon
Gibbs . . . . . . . . . . . . . . . . . . . . . . 311
Stokes . . . . . . . . . . . . . . . . . . . . . . . 38
Physical applications
electricity . . . . . . . . . . . . . . . . . . 43, 355
electromagnetism . 199, 228, 348, 358,
375, 414422, 479, 480, 484
electrostatics . . . . . . 153, 165, 194, 216,
325, 477
harmonic oscillator . . . . 239, 409413
hydrodynamics . . . . . . . . . . 5, 167, 174
mechanics . . . . . . . . . . . . . . . . . . . . 15
optics . . . . . . . . . . . . . . 38, 43, 227, 244,
314321, 324, 365, 371, 375
quantum mechanics7, 38, 128, 260, 359,
403, 377405, 427432, 497, 503
radioactivity . . . . . . . . . . . . . . . . . . 533
relativity . . . . . . . . . . . . . . 200, 211, 242
resonance . . . . . . . . . . . . . . . . . . . . . 357
thermodynamics . . 173, 175, 422427
Physical optics . . . . . . . . . . . . . 314321, 366
coherence . . . . . . . . . . . . . . . . . . 43, 371
Kirchhoff integral . . . . . . . . . 244, 349
Picard iteration . . . . . . . . . . . . . . . . . . . . . . 46
Piecewise continuous function . . . . . . . . 53
Plasma frequency . . . . . . . . . . . . . . . . . . . 228
Poincar
formula . . . . . . . . . . . . . . . . . . 516, 606
theorem . . . . . . . . . . . . . . . . . . 474, 475
Poincar Henri . . . . . . . . . . . . . . . . . 37, 475
Point
accumulation . . . . . . . . . . . . . . . 106
branch . . . . . . . . . . . . . 111, 136, 139
fixed theorem . . . . . . . . . . . . . 15, 43
saddle . . . . . . . . . . . . . . . . . . . . . . 149
stopping . . . . . . . . . . . . . . . . . . . . 167
Pointwise . . . . . . . . . . . . . . . . . . . . see Simple
Poisson
distribution . . . . . . . . . . . 530, 530, 547
equation . . . . . . . . . . . . . . 164, 216, 217
kernel . . . . . . . . . . . . . . . . . . . . . . . . 172
law . . . . . . . . . . . . . . . . . . . . . . . . . . . 616
summation formula . . . . . . . 271, 309
Poisson Denis . . . . . . . . . . . . . . . . . . . . . 172
Pole . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 110
at infinity . . . . . . . . . . . . . . . . . . . . 146
order of a . . . . . . . . . . . . . . . . . . 112
simple . . . . . . . . . . . . . . . . . . . . . . 112
Position
average . . . . . . . . . . . . . . . . . . . . . 361
operator . . . . . . . . . . 360, 391, 402, 403
representation . . . . . . . . . . . . . . . . . 360
639
uncertainty . . . . . . . . . . . . . . . . . . . 361
Positive (negative) part . . . . . . . . . . . . . . . 64
Potential
Coulomb . . . . . . . . . . . . . . . . . . . 166
Fourier transform of the . . . 306
Laplacian of the . . . . . . . . . . 208
Debye . . . . . . . . . . . . . . . . . . . . . . 325
Yukawa . . . . . . . . . . . . . . . . . . . . . 484
Power series expansion . . . . . . . . . . . . . . . 34
Pre-Hilbert space . . . . . . . . . . . . . . . . . . . 251
Principal value pv(1/x) . . . . . . . . . 188, 224
Principle
maximum . . . . . . . . . . . . . . . . . . 141
maximum modulus . . . . . . . . . . 104
uncertainty . . . . . . . . . . . . . 363, 404
Probability
conditional . . . . . . . . . . . . . . . . . 517
density . . . . . . . . . . . . . . . . . . . . . . . 526
joint . . . . . . . . . . . . . . . . 534, 539
diffuse . . . . . . . . . . . . . . . . . . . . . 525
distribution . . . . . . . . . . . . . . . . . . . 522
measure . . . . . . . . . . . . . . . . . . . . . . 514
Problem
Cauchy . . . . . . . . . . . . . . . . 238, 346
Dirichlet . . . . . . . . . . . . . . . . . . . 170
for a disc . . . . . . . . . . . . . . . . . . 172
on a half-plane . . . . . . . . . . . . . 175
for a strip . . . . . . . . . . . . . . . . . . 175
Neumann . . . . . . . . . . . . . . . . . . . 167
Proca Lagrangian . . . . . . . . . . . . . . . . . . . 484
Product
Cauchy . . . . . . . . . . . . . . . . . . . . . 215
convolution
of causal functions . . . . . . . . . . 339
of distributions . . . . . . . . . . . . . 214
and Fourier transform . . . . . . . 292
of functions . . . . . . . . . . . . . . . . 270
of functions . . . . . . . . . . . . . . . . 211
direct . . . . . . . . . see Tensor, product
exterior . . . . . . . . . . . . . . . . . 467, 468
hermitian . . . . . . . . . . . . . . . . . . . 251
scalar . . . . . . . . . . . . . . . . . . 251, 448
tensor . . . . . . . . . see Tensor product
Projection (orthogonal ) . . . . . . . 255, 399
Propagator . . . . . . . . . . . see Green function
Feynman . . . . . . . . . . . . . . . . . . . 432
Pseudo-metric . . . . . . . . . . . . . . . . . . . . . . 447
Punctured neighborhood . . . . . . . . . . . . . 18
pv(1/x) . . . . . . . . . . . . . . . . . . . . . . . 188, 224
640
R
+
R . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 59
Radius of convergence . . . . . . . . . . . . . . . 34
Ramanujan Srinivasa . . . . . . . . . . . . . . . 128
Random variable . . . . . . . . . . . 521, 521550
centered . . . . . . . . . . . . . . . . . . . . 532
discrete . . . . . . . . . . . . . . . . . . . . . 526
independent s . . . . . . . . . . . . . . . 537
product of s . . . . . . . . . . . . . . . . . 546
ratio of s . . . . . . . . . . . . . . . . . . . . 546
reduced . . . . . . . . . . . . . . . . . . . . 532
sum of s . . . . . . . . . . . . . . . . 544, 545
uncorrelated s . . . . . . . . . . . . . . . 536
Random vector . . . . . . . . . . . . . . . . . . . . . 538
Rapidly decaying function . . . . . . . . . . . 288
Rayleigh criterion . . . . . . . . . . . . . . . . . . 321
Realization of an event . . . . . . . . . . . . . . 513
Rectangle function (x) . . . . . . . . 213, 279
Regular distribution . . . . . . . . . . . . . . . . 184
Regularization of a distribution . . 234, 235
Relation
closure . . . . . . . . . . . . . . . . . . . . . 384
Kramers-Kronig dispersion . . . 227
uncertainty . . . . . . . . . 363, 364, 404
Removable singularity . . . . . . . . . . . . . . . 109
Representation
by a matrix . . . . . . . . . . . . . . . . . . . 594
faithful . . . . . . . . . . . . . . . . . . . . . 492
group . . . . . . . . . . . . . . . . . . . . . . 492
momentum . . . . . . . . . . . . . 360, 361
position . . . . . . . . . . . . . . . . . . . . 360
Riesz theorem . . . . . . . . . . . . . . . 381
trivial . . . . . . . . . . . . . . . . . . . . . . 492
Residue . . . . . . . . . . . . . . . . . . . . . . . . . . . . 114
applications . . . . . . . . . . . . . . . 117124
at infinity . . . . . . . . . . . . . . . . . . . . 146
and Fourier transform . . . . . . . . . 120
practical computation . . . . . . . . . . 116
theorem . . . . . . . . . . . . . . . . . . . . . . 115
Resolvant set . . . . . . . . . . . . . . . . . . . 393, 403
Resonance . . . . . . . . . . . . . . . . . . . . . . . . . 357
Response (impulse ) . . . . . . . . . . . . . . . 218
Riemann
integral . . . . . . . . . . . . . . . . . . . . . . . . 53
-Lebesgue lemma . . . . . . . . 267, 282
mapping theorem . . . . . . . . . . . . . 157
Riemann-Stieltjes integral . . 529, 552
sphere . . . . . . . . . . . . . . . . . . . . 146, 504
surface . . . . . . . . . . . . . . . . . . . . . . . 137
Riemann Bernhard . . . . . . . . . . . . . . . . . 138
Riesz representation theorem . . . . . . . . 381
Riesz-Fischer theorem . . . . . . . . . . . . . . . . 67
r.v. . . . . . . . . . . . . . . . . . see Random variable
Index
S
SO(3) . . . . . . . . . . . . . . . . . . . . . 492, 492505
S . . . . . . . . . . . . . . . . . . . . . . . . . . . . 289, 300
S (operator of class ) . . . . . . . . . . . . . 397
S . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 300
SU(2) . . . . . . . . . . . . . . . . . . . . . 497, 492505
S . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 26
Sn . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 465
(X ) . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 528
Saddle point . . . . . . . . . . . . . . . . . . . . . . . 149
Scalar
product . . . . . . . . . . . . . . . . . . 251, 448
pseudo-product . . . . . . . . . . . . . . . 448
Schmidt Erhard . . . . . . . . . . . . . . . . . . . 258
Schrdinger
equation . . . . . . . . . . . . . . . . . . . . . . 427
kernel . . . . . . . . . . . . . . . . . . . . . . . . 428
operator . . . . . . . . . . . . . . . . . . . . . . 427
Schwartz Laurent . . . . . . . . . . . . . . . . . . 220
Schwartz space S . . . . . . . . . . . . . . 289, 300
Schwarz Hermann . . . . . . . . . . . . . . . . . 161
Schwarz-Christoffel transformation . . . 161
Screen (total ) . . . . . . . . . . . . . . . . . . . . . 330
Self-adjoint
endomorphism . . . . . . . . . . . . . . . . 379
operator . . . . . . . . . . . . . . . . . . . . . . 394
Self-coherence function . . . . . . . . . . . . . 369
Seminorm . . . . . . . . . . . . . . . . . . . . . . . . . 577
Separable (Hilbert) space . . . . . . . . . . . . 259
Sequence (Dirac ) . . . . . . . . . . . . . . . . . 231
Series
computation of by residues . . . 122
conditionnally convergent . . . . . 26
Fourier . . . . . . . . . . . . . 258, 264, 265
partial . . . . . . . . . . . . . . . . . . . 258
Laurent . . . . . . . . . . . . . . . . . 111, 112
power . . . . . . . . . . . . . . . . . . . . . . . 34
Taylor . . . . . . . . . . . . . . . . . . . . . . . 35
Set
measurable . . . . . . . . . . . . . . . . . . 59
negligible . . . . . . . . . . . . . . . . . . . . 61
resolvant . . . . . . . . . . . . . . . 393, 403
of subsets of a set . . . . . . . . . . . . . . 57
sgn (sign distribution) . . . . . . . . . . . . 203
Si(x) . . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 312
Signal
analytic . . . . . . . . . . . . . . . . 365, 366
finite energy . . . . . . . . . . . . . . . . 368
imaginary . . . . . . . . . . . . . . . . . . 366
nonmonochromatic . . . . . . . . . 372
Signature . . . . . . . . . . . . . . . . . . . . . . . . . . 465
Simple
connectedness . . . . . . . . . . . . . . . . . 575
convergence . . . . . . . . . . . . . . . . . 19, 29
Index
curve . . . . . . . . . . . . . . . . . . . . . . . . . . 93
function . . . . . . . . . . . . . . . . . . . . . . . 63
holomorphic connectedness . . 114
pole . . . . . . . . . . . . . . . . . . . . . . . . . . 112
Simply connected . . . . . . . . . . . . . . . 114, 575
Simultaneous realization . . . . . . . . . . . . 513
Sine cardinal sinc(x) . . . . . . . . 279, 285, 288
Sine integral Si(x) . . . . . . . . . . . . . . . . . . 312
Singular distribution . . . . . . . . . . . . . . . . 185
Singularity
essential . . . . . . . . . . . . . . . . . . . . 110
at infinity . . . . . . . . . . . . . . . . . 111, 146
removable . . . . . . . . . . . . . . . . . . 109
Slowly increasing function . . . . . . . . . . . 301
Space
connected . . . . . . . . . . . . . . . . . . 574
of distributions D . . . . . . . . . . . . 183
Hilbert . . . . . . . . . . . . . . . . . . . . . 257
measurable . . . . . . . . . . . . . . . . . . 59
measure . . . . . . . . . . . . . . . . . . . . . 59
Minkowski . . . . . . . . . . . . . . . . . 448
pre-Hilbert . . . . . . . . . . . . . . . . . . 251
probability . . . . . . . . . . . . . . 512, 514
sample . . . . . . . . . . . . . . . . . . . . . 512
Schwartz S . . . . . . . . . . . . 289, 300
separable . . . . . . . . . . . . . . . . . . . 259
test D . . . . . . . . . . . . . . . . . . . . . . 182
topological . . . . . . . . . . . . . . . . . 573
Special relativity . . . . . . . . . . . . 200, 211, 242
Spectrum
continuous . . . . . . . . . . . . . 392, 403
discrete . . . . . . . . . . . . . . . . . 391, 403
of a function . . . . . . . . . . . . . . . . . 278
power . . . . . . . . . . . . . . . . . . . . . . 370
residual . . . . . . . . . . . . . . . . 393, 403
Sphere (Riemann ) . . . . . . . . . . . . . . . . 504
Spinor . . . . . . . . . . . . . . . . . . . . . . . . . . . . 503
Standard deviation . . . . . . . . . . . . . 528, 528
Star-shaped open set . . . . . . . . . . . . . . . . 475
Step (of a subdivision) . . . . . . . . . . . . . . . 52
Stieltjes (Riemann- integral) . . . 529, 552
Stirling formula . . . . . . . . . . . . . . . . . . . . 154
Stokes
formula . . . . . . . . . . . . . . . . . . . . . . 473
phenomenon . . . . . . . . . . . . . . . . . . 38
Stokes George . . . . . . . . . . . . . . . . . . . . . 472
Stopping point . . . . . . . . . . . . . . . . . . . . . 167
Structure constants . . . . . . . . . . . . . . . . . 496
Sub-vector space generated . . . . . . . . . . . 250
Subdivision . . . . . . . . . . . . . . . . . . . . . . . . . 52
Support of a distribution . . . . . . . . . . . . 187
Surface (Riemann ) . . . . . . . . . . . . . . . . 137
Symmetric
endomorphism . . . . . . . . . . . . . . . . 379
641
group . . . . . . . . . . . . . see Cyclic group
operator . . . . . . . . . . . . . . . . . . . . . . 394
System
causal . . . . . . . . . . . . . . . . . . . . . . 219
complete of events . . . . . . . . . . . 513
orthonormal . . . . . . . . . . . . . . . . 253
total . . . . . . . . . . . . . . . 257, 259, 385
T
T (C ) . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 513
Taylor
-Lagrange formula . . . . . . . . . . . . 31
-Lagrange inequality . . . . . . . . . . 32
-Young formula . . . . . . . . . . . . . . 32
expansion . . . . . . . . . . . . . . . . . . . . . 36
formula with integral remainder . 31
polynomial . . . . . . . . . . . . . . . . . . . . 31
remainder . . . . . . . . . . . . . . . . . . . . . 31
series . . . . . . . . . . . . . . . . . . . . . . . . . . 35
Taylor Brook . . . . . . . . . . . . . . . . . . . . . . . 31
Tchebychev Pafnouti . . . . . . . . . . . . . . . 548
Tchebychev inequality . . . . . . . . . . . . . . . 547
Tempered distribution . . . . . . . . . . . . . . 300
Tensor . . . . . . . . . . . . . . . . . . . . . . . . . . . . 446
Faraday . . . . . . . . . . . . . . . . . . . . . 481
Levi-Civita . . . . . . . . . . . . . . 483, 496
product . . . . . . . . . . . . . . . . . . 439447
of a linear form and a vector . 444
of linear forms . . . . . . . . . . . . . 441
of functions . . . . . . . . . . . . . . . . 209
of vector spaces . . . . . . . . . . . . . 439
of a vector and a linear form . 444
of vectors . . . . . . . . . . . . . . . . . . 443
Test function/space . . . . . . . . . . . . . . . . . 182
Theorem
Banach . . . . . . . . . . . . . . . . . . . . . . 15
Beppo Levi . . . . . . . . . . . . . . . . . . 76
Casorati-Weierstrass . . . . . . . . . . 109
Cauchy . . . . . . . . . . . . 97, 98, 98, 598
Cauchy-Lipschitz . . . . . . . . . . . . . 46
central limit . . . . . . . . . . . . 425, 558
dominated convergence . . . . . . . 75
Dini . . . . . . . . . . . . . . . . . . . . . . . . 22
Dirichlet . . . . . . . . . . . . . . . . . . . 268
Egorov . . . . . . . . . . . . . . . . . . . . . . 69
Faltung . . . . . . . . . . . . . . . . . . . . . 292
fixed point . . . . . . . . . . . . . . . . . . 15
Fubini-Lebesgue . . . . . . . . . . . . . . 79
Fubini-Tonelli . . . . . . . . . . . . . . . . 80
generalized spectral . . . . . . . . . . 398
Green . . . . . . . . . . . . . . . . . . . . . . 106
Green-Riemann . . . . . . . . . . . . . 105
Hellinger-Toeplitz . . . . . . . . . . . 394
642
Lebesgue . . . . . . . . . . . . . . . . . . . . 75
Liouville . . . . . . . . . . . . . . . . 45, 103
mean value . . . . . . . . . . . . . . . . . 141
Morera . . . . . . . . . . . . . . . . . . . . . 105
orthogonal projection . . . . . . . 255
Poincar . . . . . . . . . . . . . . . . 474, 475
representation . . . . . . . . . . . 377, 381
residue . . . . . . . . . . . . . . . . . . . . . 115
Riemann mapping . . . . . . . . . . . 157
Riesz . . . . . . . . . . . . . . . . . . . . . . . 381
Riesz-Fischer . . . . . . . . . . . . . . . . . 67
Schwarz-Christoffel . . . . . . . . . . 161
spectral . . . . . . . . . . . . . . . . . . . . . 379
Stokes . . . . . . . . . . . . . . . . . . . . . . 473
van Cittert-Zernike . . . . . . . . . . 375
Wiener-Khintchine . . . . . . . . . . 371
Total screen . . . . . . . . . . . . . . . . . . . . . . . . 330
Total system . . . . . . . . . . . . . . . 257, 259, 385
Transfer function . . . . . . . . . . . . . . . . . . . 356
Transform
z . . . . . . . . . . . . . . . . . . . . . . . . . . 344
Fourier . . . . . see Fourier, transform
Hankel . . . . . . . . . . . . . . . . . . . . . 324
Hilbert . . . . . . . . . . . . . . . . . 227, 366
Laplace . . . . . . . . . . . . . 332, 331350
inversion of the . . . . . . . . . . 338
Transformation
conformal . . . . . . . . . . . . . . . . . . 156
Schwarz-Christoffel . . . . . . . . . . 161
Translate of a distribution . . . . . . . . . . . 189
Transpose of a distribution . . . . . . . . . . 190
Transverse fields . . . . . . . . . . . . . . . . . . . . 358
Trivial representation . . . . . . . . . . . . . . . 492
U
Uncertainty
position/momentum . . . . . . . . 361
relation . . . . . . . . . . . . . . . . . . 363, 404
Uncorrelated r. v. . . . . . . . . . . . . . . . . . . . 536
Uniform convergence . . . . . . . . . . . . . . . . 29
Unitary matrix . . . . . . . . . . . . . . . . . . . . . 497
V
van Cittert-Zernike theorem . . . . . . . . . 375
Variance . . . . . . . . . . . . . . . . . . . 528, 528, 535
Vector . . . . . . . . . . . . . . . . . . . . . . . . . . . . . 497
random . . . . . . . . . . . . . . . . . . . . 538
Vector space
complete normed . . . . . . . . . . . . . 14
normed . . . . . . . . . . . . . . . . . . . . 577
sub- generated . . . . . . . . . . . . . . . 250
Index
W
Wages of fear (the ) . . . . . . . . . . . . . . . . 243
Wavelets . . . . . . . . . . . . . . . . . . . . . . . 263, 321
Weak convergence . . . . . . . . . . . . . . . . . . 230
Weierstrass Karl . . . . . . . . . . . . . . . . . . 582
Wiener-Khintchine theorem . . . . . . . . . 371
Winding number of a path . . . . . . . 96, 597
Wronskian . . . . . . . . . . . . . . . . . . . . . . . . . 409
X
X (position operator) . . . . . . . . . . . . . . . 391
Y
Young (Taylor- formula) . . . . . . . . . . . . 32
Yukawa potential . . . . . . . . . . . . . . . . . . . 484
Z
z-transform . . . . . . . . . . . . . . . . . . . . . . . . 344
Zernike (van Cittert- theorem) . . . . . 375
Zero of an holomorphic function . . . . 106
Zorns lemma . . . . . . . . . . . . . . . . . . . . . . 249

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Mathematics for Physics

Mathematics for Physics

Princeton University Press

2007 by Princeton University Press

Published by Princeton University Press, 41 William Street,

Translated from the French by Emmanuel Kowalski

All Rights Reserved

Formatted with LATEX 2 under a Linux environment.

Reminders: convergence of sequences and series

Measure theory and the Lebesgue integral

Applications to the computation of horrifying integrals or ghastly sums 117

7.2.d Other examples . . . . . . . . . . . . . . . . . . . . . . . . . . .

Hilbert spaces; Fourier series

11 Fourier transform of distributions

10 Fourier transform of functions

12 The Laplace transform

13 Physical applications of the Fourier transform

14 Bras, kets, and all that sort of thing

15.5 Quantum mechanics . . . . . . . . . . . . . . . . . . . . . . . . . . . . 427

17.6 Electromagnetism in the language of differential forms . . . . . . . . 480

19 Introduction to probability theory

20.10.b Application: Buffons needle . . . . . . . . . . . . . . . . . . . 549

A Reminders concerning topology and normed vector spaces

B Elementary reminders of differential calculus

To Kirone Mallick, Angel Alastuey and Michael Kiessling,

Remarks on the biographical summaries

Encyclopdia Universalis, Paris, 1990.

dual basis of the basis (e )

(normed) vector space

Two paradoxes involving kinetic energy

Consider a truck with mass m driving at constant speed v = 60 mph on a

Reminders concerning convergence

The problem of limits in physics

after rebounding, the speed of the ball is v2 = v2 vobs = 20 ms1 , and

Reminders concerning convergence

of conservation of total momentum).1 The kinetic energy of the system at

Hence the difference of kinetic energy is preserved, as we expected. So even

The problem of limits in physics

Fig. 1.1 Romeo, Juliet, and the boat on a lake.

Romeo, Juliet, and viscous fluids

Here is an example in mechanics where a function f (x) is defined on [0, +[,

Reminders concerning convergence

Fig. 1.2 Romeo moved closer.

or = 0. Since 6= 0, we have = 0, whichever way Romeo moved to the

The problem of limits in physics

Conclusion: The limit of a viscous fluid to a perfect fluid is singular. It is not

1.1.c Potential wall in quantum mechanics

Reminders concerning convergence

Fig. 1.4 Smoothed potential V (x) = V0 /(1 + e x/a ), with a > 0.

The particle goes through the obstacle with probability 1.

The problem of limits in physics

On the other hand, if we keep h} fixed and let a to to 0, we are converging

So the two limits [ h} 0] and [a 0] cannot be taken in an arbitrary order:

To speak of R(0, 0) has a priori no physical sense.

1.1.d Semi-infinite filter behaving as waveguide

made up of a cascading sequence of T cells (2C , L, 2C ),

Reminders concerning convergence

Hence we obtain an equation involving Z:

Some computations yield a second-degree equation, with solutions given by